Technical Documentation: modelscope/DiffSynth-Studio
ℹ️ Provenance: Multi-Tier Ingestion · CodeWiki Reference
Overview (`README.md`)
DiffSynth-Studio
Introduction
Welcome to the magical world of Diffusion models! DiffSynth-Studio is an open-source Diffusion model engine developed and maintained by the ModelScope Community team. We hope to foster technological innovation through framework construction, aggregate the power of the open-source community, and explore the rich capabilities of generative model technology!
Framework features:
- Model Support: The framework integrates mainstream open-source Diffusion models, covering image generation, video generation, audio generation, and image quality metrics models.
- VRAM Management: Dynamically schedules model parameters across disk, memory, and VRAM, allowing consumer-grade GPUs with low VRAM to run inference with large models.
- Parameter Quantization: Converts model parameters to quantized precisions such as NF4 and INT8, significantly reducing the VRAM requirements of model inference and LoRA training.
- Arbitrary Training: Almost every model that supports inference also supports training, whether base models, LoRAs, or any Adapter models with additional inputs.
- Split Training: Uses a computational graph inference engine to track every variable in the Pipeline, splitting the training process into two stages for efficient training.
References:
- Developer documentation (for humans): 中文版、English version
- Agent Skills (for AI): DiffSynth-Studio Model Integration Skills
See also:
- DiffSynth-ComfyUI: A ComfyUI node package built on DiffSynth-Studio, enabling you to build workflows in ComfyUI to run model inference.
- DiffSynth-WebUI: A lightweight LoRA training tool built on DiffSynth-Studio, enabling LoRA training of models on consumer-grade GPUs.
- ModelScope AIGC Zone (for Chinese users): Productized features powered by DiffSynth-Studio as the core inference and training engine; experience the rich potential of the open-source model ecosystem.
- ModelScope Civision (for global users): Unlock the vast potential of the open-source model ecosystem through productized capabilities powered by DiffSynth-Studio.
Update History
DiffSynth-Studio has undergone major version updates, and some old features are no longer maintained. If you need to use old features, please switch to the last historical version before the major version update.
Currently, the development personnel of this project are limited, with most of the work handled by Artiprocher and mi804. Therefore, the progress of new feature development will be relatively slow, and the speed of responding to and resolving issues is limited. We apologize for this and ask developers to understand.
October 21, 2026 We now support Fibo 2.0, including unified text-to-image and image editing inference, masked and multi-reference editing, text rendering, low VRAM inference, and full/LoRA training for both text-to-image and image editing on Fibo-2.0-dev. See the documentation and example code.
October 20, 2026 We now support MOSS-Image, including inference and training for MOSS-Image-Base and MOSS-Image-SFT. In addition, we provide an EntroPack low-bit quantized model, which can run the 3.5-bit quantized model with as little as 4 GB of VRAM; and an accelerated model distillation-accelerated via OPD training, supporting 8-step ultra-fast inference.
October 16, 2026 flux-3-image-klein open-sourced, welcome a new member to the image model family! Support includes text-to-image generation, multi-reference editing, layout control, low VRAM inference, and Base training capabilities. For details, please refer to the documentation and example code. We quantized the model with EntroPack, providing pre-quantized models in 3-bit, 4-bit, and 6-bit versions at DiffSynth-Studio/flux-3-image-klein-entropack.
October 12, 2026 We now support distillation accelerated training based on OPD (On-Policy Distillation), which enables distillation acceleration of models at an extremely low cost. For details, please refer to Distillation Accelerated Training.
September 29, 2026 We are excited to open-source EntroPack, a general-purpose tensor compression library for PyTorch that provides lossless compression with exact recovery and rate-controlled lossy compression, with GPU encoding and decoding. It is integrated into DiffSynth-Studio as a quantization backend, compressing model weights at arbitrary target bitrates in both inference and LoRA training. For more information, please refer to:
- Project Page: https://github.com/modelscope/entropack
- Technical Report: https://arxiv.org/abs/2609.34185
- Pre-quantized Models: https://www.modelscope.cn/collections/DiffSynth-Studio/EntroPack-Quantized-Models
- Example Code: Z-Image-Turbo, Qwen-Image-2.1, MiniMax-H3, LoRA Training Example
September 20, 2026 Qwen-Image-2.1 was open-sourced. The model supports text-to-image generation, image editing, and RGBA output with a transparency channel. DiffSynth-Studio now provides full support for it. For details, please refer to the documentation and example code.
September 14, 2026 We trained and open-sourced DiffSynth-Music, a controllable music generation model based on ACE-Step, which supports five control modes: Beats, Vocals, Accompany, Prosody, and Reference. For more information, please refer to:
- Model: https://modelscope.cn/models/DiffSynth-Studio/DiffSynth-Music
- Technical Report: https://arxiv.org/abs/2609.12774
- Code: ./examples/diffsynth_music/
September 11, 2026 We have added support for YuE2, a new music generation model. For details, please refer to the documentation and example code.
September 10, 2026 We have integrated LTX-2.5, the latest audio-video generation model from Lightricks. The features include text-to-audio/video with automatic duration prediction, image-to-audio/video with keyframe interpolation, audio-to-video, audio-video retake, IC-LoRA pixel spatial upscaling, text-to-audio, INT8 quantized inference, low VRAM inference, and training. For details, please refer to the documentation and code.
September 10, 2026 We have open-sourced DiffSynth-ComfyUI. Now you can run model inference with DiffSynth-Studio in ComfyUI.
September 8, 2026 We trained a MiniMax-H3 Training Adapter for the CFG-distilled MiniMax-H3 base, a DeCFG LoRA training adapter offered in FL2VA and Ref2VA versions. The training data comes from the self-generated MiniMax-H3-Self-Generated-Dataset. Based on this adapter, we trained two toy models, MiniMax-H3-Songyu-LoRA (character identity, FL2VA) and MiniMax-H3-Ref2VA-FirstFrame-Lineart (lineart first-frame control, Ref2VA), as references for inference and fine-tuning.
September 1, 2026 We have integrated SenseNova-U1.5, SenseTime's unified multimodal model, for which we provide text-to-image generation, image editing, low VRAM inference, and training support. For details, please refer to the documentation and example code.
August 31, 2026 We have integrated Qwen-Video-Edit, a video editing model developed by open-source community contributor yunpeng1998 based on the image editing model Qwen-Image-Edit. This serves as an excellent example of exploring and expanding model capabilities.
August 25, 2026 We have open-sourced DiffSynth-WebUI, enabling one-click private deployment of LoRA training services. Combined with the model quantization feature, you can train large models even with consumer-grade GPUs.
August 19, 2026 We have released the model quantization feature. It provides a unified
QuantizeConfigentry point supporting multiple quantization backends including bitsandbytes, torchao, and comfy-kitchen, with capabilities such as online quantization, loading pre-quantized weights, mixed quantization, saving quantized models, and quantization + LoRA training. For details, please refer to the documentation.
More
August 17, 2026 MiniMax-Music3 open-sourced, welcome a new member to the audio model family! Support includes text-to-music generation and low VRAM inference. For details, please refer to the documentation and example code.
August 7, 2026 We add support for Wan-Animate-2 in the Wan series. Given a reference image and a driving video, it makes the reference character perform the motions in the driving video, generating high-quality character animation, with both standard and distilled variants. For details, please refer to the documentation and example code.
August 3, 2026 MiniMax-H3 open-sourced, welcome a new member to the video model family! Support includes text-to-video-audio generation, keyframe-guided generation, reference-driven generation, low VRAM inference, and NF4-quantized inference. For details, please refer to the documentation and example code.
July 28, 2026 LingBot-Video open-sourced, welcome a new member to the video model family! This release includes two variants, Dense-1.3B and MoE-30B-A3B (30B total parameters, ~3B active per token), both supporting text-to-video, image-to-video and text-to-image generation, low VRAM inference, and LoRA / full training capabilities. For details, please refer to the documentation and example code. Huge thanks to NancyFyong for contributing the integration of this model!
July 21, 2026 We have open-sourced DiffSynth-Studio Model Integration Skills. This is a composable collection of Agent Skills that automates the entire workflow of integrating external diffusion models into DiffSynth-Studio, significantly improving the standardization and efficiency of model integration. Get started with the example!
June 29, 2026 Boogu-Image open-sourced. Support includes text-to-image generation, image editing, low VRAM inference, and training capabilities. For details, please refer to the documentation and example code.
June 24, 2026 Krea-2 is now open-source, and we have provided full support. For more details, please refer to the documentation and example code.
June 16, 2026: We have added a new Template model for ACE-Step: vocals2music. For more details, please refer to the documentation and example code.
June 15, 2026 We have open-sourced Image-to-LoRA V2, compressing the hours-long training process for image style LoRAs into a single model inference step, thereby exploring a new paradigm for LoRA model training. The technical report has been released. This release includes three models:
- DiffSynth-Studio/ZImage-i2L-v2: Adapted for the Z-Image model
- DiffSynth-Studio/KleinBase4B-i2L-v2: Adapted for the FLUX.2-klein-base-4B model
- DiffSynth-Studio/HidreamO1-i2L-v2: Adapted for the Hidream-O1-Image model
June 5, 2026 Ideogram 4 open-sourced. Support includes text-to-image inference. For details, please refer to the documentation and example code.
May 21, 2026: Added support for image quality metrics models, including FID, CLIP, Aesthetic, PickScore, ImageReward, HPSv2, and HPSv3. For details, refer to the documentation and example code.
May 18, 2026 Added CPU Offload Training support. By moving model weights layer-by-layer between CPU and GPU, it significantly reduces GPU VRAM usage during training, enabling LoRA training of large models even on consumer-grade GPUs, compatible with all models. Simply add
--enable_model_cpu_offloadto your training command to enable (currently supports single-GPU training only). For details, see the documentation.May 14, 2026 HiDream-O1-Image open-sourced, welcome a new member to the image model family! Support includes text-to-image generation, image editing, low VRAM inference, and training capabilities. For details, please refer to the documentation and example code.
April 28, 2026 We released Diffusion Templates, a plugin framework designed for Diffusion models that significantly lowers the barrier to training controllable generative models. Let's explore this cutting-edge technology together!
- Open-source code: DiffSynth-Studio
- Technical report: arXiv
- Project homepage: GitHub
- Documentation: English Version | Chinese Version
- Online demo: ModelScope
- Model collections: ModelScope | ModelScope International | HuggingFace
- Datasets: ModelScope | ModelScope International | HuggingFace
April 27, 2026 We support ACE-Step-1.5! Support includes text-to-music generation, low VRAM inference, and LoRA training capabilities. For details, please refer to the documentation and example code.
April 27, 2026: We have reinstated support for the Stable Diffusion v1.5 and SDXL models, providing academic research support exclusively for these two model types.
April 14, 2026 JoyAI-Image open-sourced, welcome a new member to the image editing model family! Support includes instruction-guided image editing, low VRAM inference, and training capabilities. For details, please refer to the documentation and example code.
March 19, 2026: Added support for openmoss/MOVA-720p and openmoss/MOVA-360p models, including training and inference capabilities. Documentation and example code are now available.
March 12, 2026: We have added support for the LTX-2.3 audio-video generation model. The features includes text-to-audio/video, image-to-audio/video, IC-LoRA control, audio-to-video, and audio-video inpainting. We have supported the complete inference and training functionalities. For details, please refer to the documentation and code.
March 3, 2026: We released the DiffSynth-Studio/Qwen-Image-Layered-Control-V2 model, which is an updated version of Qwen-Image-Layered-Control. In addition to the originally supported text-guided functionality, it adds brush-controlled layer separation capabilities.
March 2, 2026 Added support for Anima. For details, please refer to the documentation. This is an interesting anime-style image generation model. We look forward to its future updates.
February 26, 2026 Added full and lora training support for the LTX-2 audio-video generation model. See the documentation for details.
February 10, 2026 Added inference support for the LTX-2 audio-video generation model. See the documentation for details. Support for model training will be implemented in the future.
February 2, 2026 The first document of the Research Tutorial series is now available, guiding you through training a small 0.1B text-to-image model from scratch. For details, see the documentation and model. We hope DiffSynth-Studio can evolve into a more powerful training framework for Diffusion models.
January 27, 2026: Z-Image is released, and our Z-Image-i2L model is released concurrently. You can use it in ModelScope Studios. For details, see the documentation.
January 19, 2026: Added support for FLUX.2-klein-4B and FLUX.2-klein-9B models, including training and inference capabilities. Documentation and example code are now available.
January 12, 2026: We trained and open-sourced a text-guided image layer separation model (Model Link). Given an input image and a textual description, the model isolates the image layer corresponding to the described content. For more details, please refer to our blog post (Chinese version, English version).
December 24, 2025: Based on Qwen-Image-Edit-2511, we trained an In-Context Editing LoRA model (Model Link). This model takes three images as input (Image A, Image B, and Image C), and automatically analyzes the transformation from Image A to Image B, then applies the same transformation to Image C to generate Image D. For more details, please refer to our blog post (Chinese version, English version).
December 9, 2025 We release a wild model based on DiffSynth-Studio 2.0: Qwen-Image-i2L (Image-to-LoRA). This model takes an image as input and outputs a LoRA. Although this version still has significant room for improvement in terms of generalization, detail preservation, and other aspects, we are open-sourcing these models to inspire more innovative research. For more details, please refer to our blog.
December 4, 2025 DiffSynth-Studio 2.0 released! Many new features online
- Documentation online: Our documentation is still continuously being optimized and updated
- VRAM Management module upgraded, supporting layer-level disk offload, releasing both memory and VRAM simultaneously
- New model support
- Z-Image Turbo: Model, Documentation, Code
- FLUX.2-dev: Model, Documentation, Code
- Training framework upgrade
- Split Training: Supports automatically splitting the training process into two stages: data processing and training (even for training ControlNet or any other model). Computations that do not require gradient backpropagation, such as text encoding and VAE encoding, are performed during the data processing stage, while other computations are handled during the training stage. Faster speed, less VRAM requirement.
- Differential LoRA Training: This is a training technique we used in ArtAug, now available for LoRA training of any model.
- FP8 Training: FP8 can be applied to any non-training model during training, i.e., models with gradients turned off or gradients that only affect LoRA weights.
November 4, 2025 Supported the ByteDance/Video-As-Prompt-Wan2.1-14B model, which is trained based on Wan 2.1 and supports generating corresponding actions based on reference videos.
October 30, 2025 Supported the meituan-longcat/LongCat-Video model, which supports text-to-video, image-to-video, and video continuation. This model uses the Wan framework for inference and training in this project.
October 27, 2025 Supported the krea/krea-realtime-video model, adding another member to the Wan model ecosystem.
September 23, 2025 DiffSynth-Studio/Qwen-Image-EliGen-Poster released! This model was jointly developed and open-sourced by us and Taobao Experience Design Team. Built upon Qwen-Image, the model is specifically designed for e-commerce poster scenarios, supporting precise partition layout control. Please refer to our sample code.
September 9, 2025 Our training framework supports various training modes. Currently adapted for Qwen-Image, in addition to the standard SFT training mode, Direct Distill is now supported. Please refer to our sample code. This feature is experimental, and we will continue to improve it to support more comprehensive model training functions.
August 28, 2025 We support Wan2.2-S2V, an audio-driven cinematic video generation model. See ./examples/wanvideo/.
August 21, 2025 DiffSynth-Studio/Qwen-Image-EliGen-V2 released! Compared to the V1 version, the training dataset has been changed to Qwen-Image-Self-Generated-Dataset, so the generated images better conform to Qwen-Image's own image distribution and style. Please refer to our sample code.
August 21, 2025 We open-sourced the DiffSynth-Studio/Qwen-Image-In-Context-Control-Union structural control LoRA model, adopting the In Context technical route, supporting multiple categories of structural control conditions, including canny, depth, lineart, softedge, normal, and openpose. Please refer to our sample code.
August 20, 2025 We open-sourced the DiffSynth-Studio/Qwen-Image-Edit-Lowres-Fix model, improving the editing effect of Qwen-Image-Edit on low-resolution image inputs. Please refer to our sample code
August 19, 2025 Qwen-Image-Edit open-sourced, welcome a new member to the image editing model family!
August 18, 2025 We trained and open-sourced the Qwen-Image inpainting ControlNet model DiffSynth-Studio/Qwen-Image-Blockwise-ControlNet-Inpaint. The model structure adopts a lightweight design. Please refer to our sample code.
August 15, 2025 We open-sourced the Qwen-Image-Self-Generated-Dataset dataset. This is an image dataset generated using the Qwen-Image model, containing 160,000
1024 x 1024images. It includes general, English text rendering, and Chinese text rendering subsets. We provide annotations for image descriptions, entities, and structural control images for each image. Developers can use this dataset to train Qwen-Image models' ControlNet and EliGen models. We aim to promote technological development through open-sourcing!August 13, 2025 We trained and open-sourced the Qwen-Image ControlNet model DiffSynth-Studio/Qwen-Image-Blockwise-ControlNet-Depth. The model structure adopts a lightweight design. Please refer to our sample code.
August 12, 2025 We trained and open-sourced the Qwen-Image ControlNet model DiffSynth-Studio/Qwen-Image-Blockwise-ControlNet-Canny. The model structure adopts a lightweight design. Please refer to our sample code.
August 11, 2025 We open-sourced the distilled acceleration model DiffSynth-Studio/Qwen-Image-Distill-LoRA for Qwen-Image, following the same training process as DiffSynth-Studio/Qwen-Image-Distill-Full, but the model structure has been modified to LoRA, thus being better compatible with other open-source ecosystem models.
August 7, 2025 We open-sourced the entity control LoRA model DiffSynth-Studio/Qwen-Image-EliGen for Qwen-Image. Qwen-Image-EliGen can achieve entity-level controlled text-to-image generation. Technical details can be found in the paper. Training dataset: EliGenTrainSet.
August 5, 2025 We open-sourced the distilled acceleration model DiffSynth-Studio/Qwen-Image-Distill-Full for Qwen-Image, achieving approximately 5x acceleration.
August 4, 2025 Qwen-Image open-sourced, welcome a new member to the image generation model family!
August 1, 2025 FLUX.1-Krea-dev open-sourced, a text-to-image model focused on aesthetic photography. We provided comprehensive support in a timely manner, including low VRAM layer-by-layer offload, LoRA training, and full training. For more details, please refer to ./examples/flux/.
July 28, 2025 Wan 2.2 open-sourced. We provided comprehensive support in a timely manner, including low VRAM layer-by-layer offload, FP8 quantization, sequence parallelism, LoRA training, and full training. For more details, please refer to ./examples/wanvideo/.
July 11, 2025 We propose Nexus-Gen, a unified framework that combines the language reasoning capabilities of Large Language Models (LLMs) with the image generation capabilities of diffusion models. This framework supports seamless image understanding, generation, and editing tasks.
- Paper: Nexus-Gen: Unified Image Understanding, Generation, and Editing via Prefilled Autoregression in Shared Embedding Space
- GitHub Repository: https://github.com/modelscope/Nexus-Gen
- Model: ModelScope, HuggingFace
- Training Dataset: ModelScope Dataset
- Online Experience: ModelScope Nexus-Gen Studio
June 15, 2025 ModelScope's official evaluation framework EvalScope now supports text-to-image generation evaluation. Please refer to the best practices guide to try it out.
March 25, 2025 Our new open-source project DiffSynth-Engine is now open-sourced! Focused on stable model deployment, targeting industry, providing better engineering support, higher computational performance, and more stable features.
March 31, 2025 We support InfiniteYou, a face feature preservation method for FLUX. More details can be found in ./examples/InfiniteYou/.
March 13, 2025 We support HunyuanVideo-I2V, the image-to-video generation version of Tencent's open-source HunyuanVideo. More details can be found in ./examples/HunyuanVideo/.
February 25, 2025 We support Wan-Video, a series of state-of-the-art video synthesis models open-sourced by Alibaba. See ./examples/wanvideo/.
February 17, 2025 We support StepVideo! Advanced video synthesis model! See ./examples/stepvideo.
December 31, 2024 We propose EliGen, a new framework for entity-level controlled text-to-image generation, supplemented with an inpainting fusion pipeline, extending its capabilities to image inpainting tasks. EliGen can seamlessly integrate existing community models such as IP-Adapter and In-Context LoRA, enhancing their versatility. For more details, see ./examples/EntityControl.
- Paper: EliGen: Entity-Level Controlled Image Generation with Regional Attention
- Model: ModelScope, HuggingFace
- Online Experience: ModelScope EliGen Studio
- Training Dataset: EliGen Train Set
December 19, 2024 We implemented advanced VRAM management for HunyuanVideo, enabling video generation with resolutions of 129x720x1280 on 24GB VRAM or 129x512x384 on just 6GB VRAM. More details can be found in ./examples/HunyuanVideo/.
December 18, 2024 We propose ArtAug, a method to improve text-to-image models through synthesis-understanding interaction. We trained an ArtAug enhancement module for FLUX.1-dev in LoRA format. This model incorporates the aesthetic understanding of Qwen2-VL-72B into FLUX.1-dev, thereby improving the quality of generated images.
- Paper: https://arxiv.org/abs/2412.12888
- Example: https://github.com/modelscope/DiffSynth-Studio/tree/main/examples/ArtAug
- Model: ModelScope, HuggingFace
- Demo: ModelScope, HuggingFace (coming soon)
October 25, 2024 We provide extensive FLUX ControlNet support. This project supports many different ControlNet models and can be freely combined, even if their structures are different. Additionally, ControlNet models are compatible with high-resolution optimization and partition control technologies, enabling very powerful controllable image generation. See
./examples/ControlNet/.October 8, 2024 We released extended LoRAs based on CogVideoX-5B and ExVideo. You can download this model from ModelScope or HuggingFace.
August 22, 2024 This project now supports CogVideoX-5B. See here. We provide several interesting features for this text-to-video model, including:
- Text-to-video
- Video editing
- Self super-resolution
- Video interpolation
August 22, 2024 We implemented an interesting brush feature that supports all text-to-image models. Now you can create stunning images with the assistance of AI using the brush!
- Use it in our WebUI.
August 21, 2024 DiffSynth-Studio now supports FLUX.
- Enable CFG and high-resolution inpainting to improve visual quality. See here
- LoRA, ControlNet, and other addon models will be released soon.
June 21, 2024 We propose ExVideo, a post-training fine-tuning technique aimed at enhancing the capabilities of video generation models. We extended Stable Video Diffusion to achieve long video generation of up to 128 frames.
- Project Page
- Source code has been released in this repository. See
examples/ExVideo. - Model has been released at HuggingFace and ModelScope.
- Technical report has been released at arXiv.
- You can try ExVideo in this demo!
June 13, 2024 DiffSynth Studio has migrated to ModelScope. The development team has also transitioned from "me" to "us". Of course, I will still participate in subsequent development and maintenance work.
January 29, 2024 We propose Diffutoon, an excellent cartoon coloring solution.
- Project Page
- Source code has been released in this project.
- Technical report (IJCAI 2024) has been released at arXiv.
December 8, 2023 We decided to initiate a new project aimed at unleashing the potential of diffusion models, especially in video synthesis. The development work of this project officially began.
November 15, 2023 We propose FastBlend, a powerful video deflickering algorithm.
October 1, 2023 We released an early version of the project named FastSDXL. This was an initial attempt to build a diffusion engine.
- Source code has been released at GitHub.
- FastSDXL includes a trainable OLSS scheduler to improve efficiency.
August 29, 2023 We propose DiffSynth, a video synthesis framework.
- Project Page.
- Source code has been released at EasyNLP.
- Technical report (ECML PKDD 2024) has been released at arXiv.
Installation
Install from source (recommended):
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more installation methods and instructions for non-NVIDIA GPUs, please refer to the Installation Guide.
Download Source Configuration
Before model inference and training, you can configure the model download source and other options through environment variables.
This project downloads models from ModelScope by default. For users outside China, you can download models from the ModelScope International site with the following configuration:
shellexport MODELSCOPE_ENDPOINT=https://modelscope.aiTo download models from HuggingFace, please modify the environment variables. Note that model IDs may differ across platforms:
shellexport DIFFSYNTH_DOWNLOAD_SOURCE="huggingface"
Basic Framework
DiffSynth-Studio redesigns the inference and training pipelines for mainstream Diffusion models (including FLUX, Wan, etc.), enabling efficient memory management and flexible model training.
Quick start: experience popular and the latest models:
| Architecture | Model ID | Inference | Low VRAM Inference | Full Training | Validation After Full Training | LoRA Training | Validation After LoRA Training |
|---|---|---|---|---|---|---|---|
| Fibo 2.0 | briaai/Fibo-2.0-dev | code | code | code | code | code | code |
| Fibo 2.0 | briaai/Fibo-2.0-edit-turbo | code | code | - | - | - | - |
| Qwen-Image-2.1 | Qwen/Qwen-Image-2.1 | code | code | code | code | code | code |
| Qwen-Image-2.1 | Qwen/Qwen-Image-2.1-Turbo | code | code | code | code | code | code |
| MiniMax-H3 | MiniMax/MiniMax-H3: FL2VA | code | code | code | code | code | code |
| MiniMax-H3 | DiffSynth-Studio/MiniMax-H3-NF4: FL2VA pruned | code | code | - | - | code | code |
| ACE-Step | ACE-Step/acestep-v15-xl-sft | code | code | code | code | code | code |
| Z-Image | Tongyi-MAI/Z-Image-Turbo | code | code | code | code | code | code |
| Krea-2 | krea/Krea-2-Raw | code | code | code | code | code | code |
| Krea-2 | krea/Krea-2-Turbo | code | code | code | code | code | code |
Model overview:
- Image generation
- Fibo 2.0: Documentation, Example code
- MOSS-Image: Documentation, Example code
- FLUX.3: Documentation, Example code
- Qwen-Image-2.1: Documentation, Example code
- SenseNova-U1: Documentation, Example code
- Boogu-Image: Documentation, Example code
- Krea-2: Documentation, Example code
- Ideogram 4: Documentation, Example code
- HiDream-O1-Image: Documentation, Example code
- JoyAI-Image: Documentation, Example code
- ERNIE-Image: Documentation, Example code
- FLUX.2: Documentation, Example code
- Z-Image: Documentation, Example code
- Anima: Documentation, Example code
- Qwen-Image: Documentation, Example code
- FLUX.1: Documentation, Example code
- Stable Diffusion XL: Documentation, Example code
- Stable Diffusion: Documentation, Example code
- Video generation
- MiniMax-H3: Documentation, Example code
- LingBot-Video: Documentation, Example code
- LTX-2: Documentation, Example code
- Wan: Documentation, Example code
- Audio generation
- DiffSynth-Music: Documentation, Example code
- MiniMax-Music3: Documentation, Example code
- YuE2: Documentation, Example code
- ACE-Step: Documentation, Example code
- Evaluation models: Documentation, Example code
Innovative Achievements
We believe that a well-developed open-source code framework can lower the threshold for technical exploration. We have achieved many interesting technologies based on this codebase. Perhaps you also have many wild ideas, and with DiffSynth-Studio, you can quickly realize these ideas.
EntroPack: A General-Purpose Tensor Compression Library for PyTorch
Compress model weights at arbitrary target bitrates with GPU encoding and decoding, and run the compressed weights directly in inference and LoRA training.
- Project Page: GitHub
- Technical Report: https://arxiv.org/abs/2609.34185
- Pre-quantized Models: ModelScope
- Documentation: English Version | Chinese Version
- Example Code: Z-Image-Turbo, Qwen-Image-2.1, MiniMax-H3, LoRA Training Example
TreeAdapter: A Model System Built from Structured LoRAs
A model system composed of 10,000+ LoRAs challenges the precise generation of tens of thousands of rare species.
- Paper: TreeAdapter: Hierarchical Taxonomy-Guided Adapter Composition for Fine-Grained Species Image Generation
- Model: ModelScope
Image-to-LoRA: Compressing Model Training into Model Inference
True Meta Learning: feed a dataset of images into one end of the model, and the trained LoRA model comes out the other end.
- Paper: Compressing Image Style Training into a Single Model Forward
- Model:
- DiffSynth-Studio/ZImage-i2L-v2: Adapted for the Z-Image model
- DiffSynth-Studio/KleinBase4B-i2L-v2: Adapted for the FLUX.2-klein-base-4B model
- DiffSynth-Studio/HidreamO1-i2L-v2: Adapted for the Hidream-O1-Image model
| Input example 1 | Output example 1 | Input example 2 | Output example 2 |
|---|---|---|---|
Diffusion-Templates: A Plugin-Based Controllable Generation Framework
One framework that turns every controllable generation capability into a plugin, allowing multiple models to combine and emerge with rich generation capabilities.
- Paper: Diffusion Templates: A Unified Plugin Framework for Controllable Diffusion
- Project homepage: GitHub
- Documentation: English Version | Chinese Version
- Online demo: ModelScope
- Model collections: ModelScope | ModelScope International | HuggingFace
- Datasets: ModelScope | ModelScope International | HuggingFace
| Reference image | Local editing | Style transfer | Sharpness enhancement |
|---|---|---|---|
Spectral Evolution Search: Efficient Inference-Time Scaling for Reward-Aligned Image Generation
Trade inference time for higher quality of generated content.
- Paper: Spectral Evolution Search: Efficient Inference-Time Scaling for Reward-Aligned Image Generation
- Sample Code: /docs/en/Research_Tutorial/inference_time_scaling.md
| FLUX.1-dev | FLUX.1-dev + SES | Qwen-Image | Qwen-Image + SES |
|---|---|---|---|
VIRAL: Visual In-Context Reasoning via Analogy in Diffusion Transformers
Transform image 3 into image 4 based on the change from image 1 to image 2 — the emergent capability of image editing models.
- Paper: VIRAL: Visual In-Context Reasoning via Analogy in Diffusion Transformers
- Sample code: /examples/qwen_image/model_inference/Qwen-Image-Edit-2511-ICEdit.py
- Model: ModelScope
| Example 1 | Example 2 | Query | Output |
|---|---|---|---|
AttriCtrl: Attribute Intensity Control for Image Generation Models
Numerical attributes can also precisely control image generation models.
- Paper: AttriCtrl: Fine-Grained Control of Aesthetic Attribute Intensity in Diffusion Models
- Sample Code: /examples/flux/model_inference/FLUX.1-dev-AttriCtrl.py
- Model: ModelScope
| brightness scale = 0.1 | brightness scale = 0.3 | brightness scale = 0.5 | brightness scale = 0.7 | brightness scale = 0.9 |
|---|---|---|---|---|
AutoLoRA: Automated LoRA Retrieval and Fusion
LoRA is a product that unifies needs and solutions — how can we make better use of these LoRAs?
- Paper: AutoLoRA: Automatic LoRA Retrieval and Fine-Grained Gated Fusion for Text-to-Image Generation
- Sample Code: /examples/flux/model_inference/FLUX.1-dev-LoRA-Fusion.py
- Model: ModelScope
| LoRA 1 | LoRA 2 | LoRA 3 | LoRA 4 | |
|---|---|---|---|---|
| LoRA 1 | ||||
| LoRA 2 | ||||
| LoRA 3 | ||||
| LoRA 4 |
Nexus-Gen: Unified Architecture for Image Understanding, Generation, and Editing
What happens when a single model combines image understanding, generation, and editing capabilities?
- Detailed Page: https://github.com/modelscope/Nexus-Gen
- Paper: Nexus-Gen: Unified Image Understanding, Generation, and Editing via Prefilled Autoregression in Shared Embedding Space
- Model: ModelScope, HuggingFace
- Dataset: ModelScope Dataset
- Online Experience: ModelScope Nexus-Gen Studio
ArtAug: Aesthetic Enhancement for Image Generation Models
A single LoRA that significantly enhances detail and aesthetics.
- Detailed Page: ./examples/ArtAug/
- Paper: ArtAug: Enhancing Text-to-Image Generation through Synthesis-Understanding Interaction
- Model: ModelScope, HuggingFace
- Online Experience: ModelScope AIGC Tab
| FLUX.1-dev | FLUX.1-dev + ArtAug LoRA |
|---|---|
EliGen: Precise Image Partition Control
How can region-based layers control the position of content in an image?
- Paper: EliGen: Entity-Level Controlled Image Generation with Regional Attention
- Sample Code: /examples/flux/model_inference/FLUX.1-dev-EliGen.py
- Model: ModelScope, HuggingFace
- Online Experience: ModelScope EliGen Studio
- Dataset: EliGen Train Set
| Entity Control Region | Generated Image |
|---|---|
ExVideo: Extended Training for Video Generation Models
If a video generation model can only generate 25 frames, how can we make it generate longer videos?
- Project Page: Project Page
- Paper: ExVideo: Extending Video Diffusion Models via Parameter-Efficient Post-Tuning
- Sample Code: Please refer to the older version
- Model: ModelScope, HuggingFace
https://github.com/modelscope/DiffSynth-Studio/assets/35051019/d97f6aa9-8064-4b5b-9d49-ed6001bb9acc
Diffutoon: High-Resolution Anime-Style Video Rendering
I don't care what you say, I just love anime!
- Project Page: Project Page
- Paper: Diffutoon: High-Resolution Editable Toon Shading via Diffusion Models
- Sample Code: Please refer to the older version
The Diffutoon examples are from the pre-2.0 codebase and are not included in the current main branch. Use the older source tree above with the corresponding older version of DiffSynth-Studio.
https://github.com/Artiprocher/DiffSynth-Studio/assets/35051019/b54c05c5-d747-4709-be5e-b39af82404dd
DiffSynth: The Original Version of This Project
In the era before video generation models, how could image generation models be used to process videos?
- Project Page: Project Page
- Paper: DiffSynth: Latent In-Iteration Deflickering for Realistic Video Synthesis
- Sample Code: Please refer to the older version
https://github.com/Artiprocher/DiffSynth-Studio/assets/35051019/59fb2f7b-8de0-4481-b79f-0c3a7361a1ea
Contact Us
| Discord:https://discord.gg/Mm9suEeUDc |
|---|
All Supported Models
Overview (`docs/en/README.md`)
DiffSynth-Studio Documentation
Welcome to the magical world of Diffusion models! DiffSynth-Studio is an open-source Diffusion model engine developed and maintained by the ModelScope Community. We aim to build a universal Diffusion model framework that fosters technological innovation through framework construction, aggregates the power of the open-source community, and explores the boundaries of generative model technology!
Documentation Reading Guide
graph LR;
use["I want to use models for inference and training"]-->sec1["Section 1: Getting Started"];
use["I want to use models for inference and training"]-->sec2["Section 2: Model Details"];
use["I want to use models for inference and training"]-->sec3["Section 3: Training Framework"];
develop["I want to develop based on this framework"]-->sec3["Section 3: Training Framework"];
develop["I want to develop based on this framework"]-->sec4["Section 4: Model Integration"];
develop["I want to develop based on this framework"]-->sec5["Section 5: API Reference"];
explore["I want to explore new technologies based on this project"]-->sec4["Section 4: Model Integration"];
explore["I want to explore new technologies based on this project"]-->sec5["Section 5: API Reference"];
explore["I want to explore new technologies based on this project"]-->sec6["Section 6: Diffusion Templates"];
explore["I want to explore new technologies based on this project"]-->sec7["Section 7: Research Guide"];
problem["I encountered a problem"]-->sec8["Section 8: Frequently Asked Questions"];Section 1: Getting Started
This section introduces the basic usage of DiffSynth-Studio, including how to enable VRAM management for inference on GPUs with extremely low VRAM, and how to train various base models, LoRAs, ControlNets, and other models.
- Installation Dependencies
- Model Inference
- Accelerated Inference
- VRAM Management
- Model Quantization
- Model Training
- Environment Variables
- GPU/NPU Support
- Inference WebUI
Section 2: Model Details
This section introduces the Diffusion models supported by DiffSynth-Studio. Some model pipelines feature special functionalities such as controllable generation and parallel acceleration.
- FLUX.1
- Wan
- Qwen-Image
- Qwen-Video-Edit
- FLUX.2
- Z-Image
- Anima
- LTX-2
- ERNIE-Image
- JoyAI-Image
- ACE-Step
- HiDream-O1-Image
- Stable Diffusion
- Stable Diffusion XL
- Image Quality Metrics
- Ideogram 4
- Krea-2
- Fibo 2.0
- Boogu-Image
- LingBot-Video
- MiniMax-H3
- MiniMax-Music3
- SenseNova-U1
- YuE2
- Qwen-Image-2.1
- FLUX.3
Section 3: Training Framework
This section introduces the design philosophy of the training framework in DiffSynth-Studio, helping developers understand the principles of Diffusion model training algorithms.
- Basic Principles of Diffusion Models
- Standard Supervised Training
- Enabling FP8 Precision in Training
- End-to-End Distillation Accelerated Training
- Two-Stage Split Training
- Differential LoRA Training
- Enabling DeepSpeed
- Offload Training
Section 4: Model Integration
This section introduces how to integrate models into DiffSynth-Studio to utilize the framework's basic functions, helping developers provide support for new models in this project or perform inference and training of private models.
- Integrating Model Architecture
- Building a Pipeline
- Enabling Fine-Grained VRAM Management
- Model Training Integration
- Integrating a Quantization Backend
We have open-sourced DiffSynth-Studio Model Integration Skills. This is a composable collection of Agent Skills that automates the entire workflow of integrating external diffusion models (image / video / audio) into DiffSynth-Studio. It distills and defines the model integration standards of DiffSynth-Studio, consolidating best practices such as codebase analysis, model code integration, Pipeline design, low-VRAM management, and training support into reusable standard procedures. Following these standards can significantly lower the integration barrier, reduce repetitive debugging, and greatly improve the efficiency of integrating new models. We recommend starting from the
diffsynth-integratorexample to get started quickly and accelerate model integration.
Section 5: API Reference
This section introduces the independent core module diffsynth.core in DiffSynth-Studio, explaining how internal functions are designed and operate. Developers can use these functional modules in other codebase developments if needed.
diffsynth.core.attention: Attention mechanism implementationdiffsynth.core.data: Data processing operators and general datasetsdiffsynth.core.gradient: Gradient checkpointingdiffsynth.core.loader: Model download and loadingdiffsynth.core.quant: Model quantizationdiffsynth.core.vram: VRAM management
Section 6: Diffusion Templates
This section introduces the controllable generation plugin framework for Diffusion models, explaining the framework's operation mechanism and how to use Template models for inference and training.
- Introducing Diffusion Templates
- Diffusion Templates Architecture Details
- Template Model Inference
- Template Model Training
Section 7: Research Guide
This section introduces how to use DiffSynth-Studio to train new models, helping researchers explore new model technologies.
- Training models from scratch
- Inference improvement techniques
- Designing controllable generation models
- Creating new training paradigms 【coming soon】
Section 8: Frequently Asked Questions
This section summarizes common developer questions. If you encounter issues during usage or development, please refer to this section. If you still cannot resolve the problem, please submit an issue on GitHub.
Overview (`docs/zh/README.md`)
DiffSynth-Studio 文档
欢迎来到 Diffusion 模型的魔法世界!DiffSynth-Studio 是由魔搭社区团队开发和维护的开源 Diffusion 模型引擎。我们期望构建一个通用的 Diffusion 模型框架,以框架建设孵化技术创新,凝聚开源社区的力量,探索生成式模型技术的边界!
文档阅读导引
graph LR;
我想要使用模型进行推理和训练-->sec1[Section 1: 上手使用];
我想要使用模型进行推理和训练-->sec2[Section 2: 模型详解];
我想要使用模型进行推理和训练-->sec3[Section 3: 训练框架];
我想要基于此框架进行二次开发-->sec3[Section 3: 训练框架];
我想要基于此框架进行二次开发-->sec4[Section 4: 模型接入];
我想要基于此框架进行二次开发-->sec5[Section 5: API 参考];
我想要基于本项目探索新的技术-->sec4[Section 4: 模型接入];
我想要基于本项目探索新的技术-->sec5[Section 5: API 参考];
我想要基于本项目探索新的技术-->sec6[Section 6: Diffusion Templates]
我想要基于本项目探索新的技术-->sec7[Section 7: 学术导引];
我遇到了问题-->sec8[Section 8: 常见问题];Section 1: 上手使用
本节介绍 DiffSynth-Studio 的基本使用方式,包括如何启用显存管理从而在极低显存的 GPU 上进行推理,以及如何训练任意基础模型、LoRA、ControlNet 等模型。
Section 2: 模型详解
本节介绍 DiffSynth-Studio 所支持的 Diffusion 模型,部分模型 Pipeline 具备可控生成、并行加速等特色功能。
- FLUX.1
- Wan
- Qwen-Image
- Qwen-Video-Edit
- FLUX.2
- Z-Image
- Anima
- LTX-2
- ERNIE-Image
- JoyAI-Image
- ACE-Step
- HiDream-O1-Image
- Stable Diffusion
- Stable Diffusion XL
- 图像质量评估指标
- Ideogram 4
- Krea-2
- Fibo 2.0
- Boogu-Image
- LingBot-Video
- MiniMax-H3
- MiniMax-Music3
- SenseNova-U1
- YuE2
- Qwen-Image-2.1
- FLUX.3
Section 3: 训练框架
本节介绍 DiffSynth-Studio 中训练框架的设计思路,帮助开发者理解 Diffusion 模型训练算法的原理。
Section 4: 模型接入
本节介绍如何将模型接入 DiffSynth-Studio 从而使用框架基础功能,帮助开发者为本项目提供新模型的支持,或进行私有化模型的推理和训练。
我们开源了 DiffSynth-Studio Model Integration Skills。这是一套可组合的 Agent Skill 合集,将外部扩散模型(图像 / 视频 / 音频)接入 DiffSynth-Studio 的全流程自动化。它沉淀并定义了 DiffSynth-Studio 的模型接入规范,将代码库分析、模型代码接入、Pipeline 设计、低显存管理与训练支持等最佳实践固化为可复用的标准流程。遵循这套规范,能显著降低接入门槛、减少反复调试,大幅提升新模型的接入效率。建议从 diffsynth-integrator 的使用示例开始,快速上手、加速模型接入。
Section 5: API 参考
本节介绍 DiffSynth-Studio 中的独立核心模块 diffsynth.core,介绍内部的功能是如何设计和运作的,开发者如有需要,可将其中的功能模块用于其他代码库的开发中。
diffsynth.core.attention: 注意力机制实现diffsynth.core.data: 数据处理算子与通用数据集diffsynth.core.gradient: 梯度检查点diffsynth.core.loader: 模型下载与加载diffsynth.core.quant: 模型量化diffsynth.core.vram: 显存管理
Section 6: Diffusion Templates
本节介绍 Diffusion 模型可控生成插件框架 Diffusion Templates,讲解 Diffusion Templates 框架的运行机制,展示如何使用 Template 模型进行推理和训练。
Section 7: 学术导引
本节介绍如何利用 DiffSynth-Studio 训练新的模型,帮助科研工作者探索新的模型技术。
Section 8: 常见问题
本节总结了开发者常见的问题,如果你在使用和开发中遇到了问题,请参考本节内容,如果仍无法解决,请到 GitHub 上给我们提 issue。
Overview (`examples/anima/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/Anima.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/Anima.html
Overview (`examples/fibo/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/Fibo.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/Fibo.html
Inference
- Text-to-image, editing, masked editing, text rendering:
model_inference/Fibo-2.0-dev.py - Editing and masked editing:
model_inference/Fibo-2.0-edit-turbo.py - Low VRAM versions:
model_inference_low_vram/
Training
Only Fibo-2.0-dev is trainable.
| Mode | Full Training | LoRA Training | Full Validation | LoRA Validation |
|---|---|---|---|---|
| Text-to-image | full/Fibo-2.0-dev-T2I.sh |
lora/Fibo-2.0-dev-T2I.sh |
validate_full/Fibo-2.0-dev-T2I.py |
validate_lora/Fibo-2.0-dev-T2I.py |
| Image editing | full/Fibo-2.0-dev-Edit.sh |
lora/Fibo-2.0-dev-Edit.sh |
validate_full/Fibo-2.0-dev-Edit.py |
validate_lora/Fibo-2.0-dev-Edit.py |
Overview (`examples/flux/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/FLUX.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/FLUX.html
Overview (`examples/flux2/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/FLUX2.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/FLUX2.html
Overview (`examples/lingbot_video/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/LingBot-Video.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/LingBot-Video.html
Overview (`examples/ltx2/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/LTX-2.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/LTX-2.html
Overview (`examples/mova/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/Wan.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/Wan.html
Overview (`examples/qwen_image/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/Qwen-Image.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/Qwen-Image.html
Overview (`examples/wanvideo/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/Wan.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/Wan.html
Overview (`examples/yue2/README.md`)
YuE2
YuE2 is a music generation model with editable symbolic planning.
- Inference: YuE2.py
- Low VRAM inference: YuE2.py
- Plan and render an edited score: YuE2-Edit.py
- Low VRAM plan and render an edited score: YuE2-Edit.py
- Documentation: English, Chinese
Overview (`examples/z_image/README.md`)
English Document: https://diffsynth-studio-doc.readthedocs.io/en/latest/Model_Details/Z-Image.html
中文文档:https://diffsynth-studio-doc.readthedocs.io/zh-cn/latest/Model_Details/Z-Image.html
Overview (`diffsynth/utils/ses/README.md`)
Please see docs/en/Research_Tutorial/inference_time_scaling.md or docs/zh/Research_Tutorial/inference_time_scaling.md for more details.
Requirements (`docs/requirements.txt`)
docutils>=0.16.0
myst_parser
recommonmark
sphinx>=5.3.0
sphinx-book-theme
sphinx-copybutton
sphinx-autobuild
sphinx-rtd-theme
sphinx_markdown_tables
sphinxcontrib-mermaid
pymdown-extensions
En .Readthedocs.Yaml (`docs/en/.readthedocs.yaml`)
.readthedocs.yaml
Read the Docs configuration file
See https://docs.readthedocs.io/en/stable/config-file/v2.html for details
Required
version: 2
Set the OS, Python version and other tools you might need
build:
os: ubuntu-22.04
tools:
python: "3.10"
Build documentation in the "docs/" directory with Sphinx
sphinx:
configuration: docs/en/conf.py
Optionally build your docs in additional formats such as PDF and ePub
formats:
- epub
Optional but recommended, declare the Python requirements required
to build your documentation
See https://docs.readthedocs.io/en/stable/guides/reproducible-builds.html
python:
install:
- requirements: docs/requirements.txt
En QA (`docs/en/QA.md`)
Frequently Asked Questions
Why doesn't the training framework support batch size > 1?
- Larger batch sizes no longer achieve significant acceleration: Due to acceleration technologies such as flash attention that have fully improved GPU utilization, larger batch sizes will only bring greater VRAM usage without significant acceleration. The experience with small models like Stable Diffusion 1.5 is no longer applicable to the latest large models.
- Larger batch sizes can be achieved through other solutions: Multi-GPU training and Gradient Accumulation can both mathematically equivalently achieve larger batch sizes.
- Larger batch sizes contradict the framework's general design: We hope to build a general training framework. Many models cannot accommodate larger batch sizes, such as text encodings of different lengths and images of different resolutions, which cannot be merged into larger batches.
Why aren't redundant parameters removed from certain models?
In some models, redundant parameters exist. For example, in Qwen-Image's DiT model, the text portion of the last layer does not participate in any calculations. This is a minor bug left by the model developers. Setting it as trainable directly will also cause errors in multi-GPU training.
To maintain compatibility with other models in the open-source community, we have decided to retain these parameters. These redundant parameters can avoid errors in multi-GPU training through the --find_unused_parameters parameter.
Why does FP8 quantization show no acceleration effect?
Native FP8 computation relies on Hopper architecture GPUs and has significant precision errors. It is currently immature technology, so this project does not support native FP8 computation.
FP8 computation in VRAM management refers to storing model parameters in memory or VRAM with FP8 precision and temporarily converting them to other precisions when needed for computation. Therefore, it can only reduce VRAM usage without acceleration effects.
Why doesn't the training framework support native FP8 precision training?
Even with suitable hardware conditions, we currently have no plans to support native FP8 precision training.
- The main challenge of native FP8 precision training is precision overflow caused by gradient explosion. To ensure training stability, the model structure needs to be redesigned accordingly. However, no model developers are willing to do so at present.
- Additionally, models trained with native FP8 precision can only be computed with BF16 precision during inference without Hopper architecture GPUs, theoretically resulting in generation quality inferior to FP8.
Therefore, native FP8 precision training technology is extremely immature. We will observe the technological developments in the open-source community.
How to dynamically load LoRA models during inference?
We support two loading methods for LoRA models. See LoRA Loading for details:
- Cold Loading: When VRAM Management is not enabled for the base model, LoRA will be fused into the base model weights. In this case, inference speed remains unchanged, and LoRA cannot be unloaded after loading.
- Hot Loading: When VRAM Management is enabled for the base model, LoRA will not be fused into the base model weights. In this case, inference speed will slow down, and LoRA can be unloaded after loading via
pipe.clear_lora().
How to reduce VRAM usage during training?
The framework supports multiple methods to reduce VRAM usage during training, including Gradient Checkpointing, DeepSpeed, FP8, Two-Stage Split Training, and CPU Offload. Please refer to Low VRAM Training.
En Index (`docs/en/index.rst`)
Welcome to DiffSynth-Studio's Documentation
.. toctree::
:maxdepth: 2
:caption: Documentation Introduction
README
.. toctree::
:maxdepth: 2
:caption: Getting Started
Pipeline_Usage/Setup
Pipeline_Usage/Model_Inference
Pipeline_Usage/Accelerated_Inference
Pipeline_Usage/VRAM_management
Pipeline_Usage/Quantization
Pipeline_Usage/Model_Training
Pipeline_Usage/Environment_Variables
Pipeline_Usage/GPU_support
Pipeline_Usage/Inference_WebUI
.. toctree::
:maxdepth: 2
:caption: Model Details
Model_Details/FLUX
Model_Details/Wan
Model_Details/Qwen-Image
Model_Details/Qwen-Video-Edit
Model_Details/FLUX2
Model_Details/Z-Image
Model_Details/Anima
Model_Details/LTX-2
Model_Details/LTX-2.5
Model_Details/ERNIE-Image
Model_Details/JoyAI-Image
Model_Details/ACE-Step
Model_Details/HiDream-O1-Image
Model_Details/Stable-Diffusion
Model_Details/Stable-Diffusion-XL
Model_Details/Image-Quality-Metrics
Model_Details/Ideogram-4
Model_Details/Krea-2
Model_Details/Boogu-Image
Model_Details/LingBot-Video
Model_Details/MiniMax-H3
Model_Details/DiffSynth-Music
Model_Details/MiniMax-Music3
Model_Details/YuE2
Model_Details/SenseNova-U1
Model_Details/Qwen-Image-2.1
Model_Details/FLUX.3
Model_Details/MOSS-Image
Model_Details/Fibo
.. toctree::
:maxdepth: 2
:caption: Training Framework
Training/Understanding_Diffusion_models
Training/Supervised_Fine_Tuning
Training/FP8_Precision
Training/Distillation
Training/Direct_Distill
Training/Split_Training
Training/Differential_LoRA
Training/DeepSpeed
Training/Offload_Training
.. toctree::
:maxdepth: 2
:caption: Model Integration
Developer_Guide/Integrating_Your_Model
Developer_Guide/Building_a_Pipeline
Developer_Guide/Enabling_VRAM_management
Developer_Guide/Training_Diffusion_Models
Developer_Guide/Integrating_Quantization_Backend
.. toctree::
:maxdepth: 2
:caption: API Reference
API_Reference/core/attention
API_Reference/core/data
API_Reference/core/gradient
API_Reference/core/loader
API_Reference/core/quant
API_Reference/core/vram
.. toctree::
:maxdepth: 2
:caption: Diffusion Templates
Diffusion_Templates/Introducing_Diffusion_Templates.md
Diffusion_Templates/Understanding_Diffusion_Templates.md
Diffusion_Templates/Template_Model_Inference.md
Diffusion_Templates/Template_Model_Training.md
.. toctree::
:maxdepth: 2
:caption: Research Guide
Research_Tutorial/train_from_scratch
Research_Tutorial/inference_time_scaling
Research_Tutorial/controllable_models
.. toctree::
:maxdepth: 2
:caption: FAQ
QA
Indices and tables
- :ref:
genindex - :ref:
modindex - :ref:
search
Zh .Readthedocs.Yaml (`docs/zh/.readthedocs.yaml`)
.readthedocs.yaml
Read the Docs configuration file
See https://docs.readthedocs.io/en/stable/config-file/v2.html for details
Required
version: 2
Set the OS, Python version and other tools you might need
build:
os: ubuntu-22.04
tools:
python: "3.10"
Build documentation in the "docs/" directory with Sphinx
sphinx:
configuration: docs/zh/conf.py
Optionally build your docs in additional formats such as PDF and ePub
formats:
- epub
Optional but recommended, declare the Python requirements required
to build your documentation
See https://docs.readthedocs.io/en/stable/guides/reproducible-builds.html
python:
install:
- requirements: docs/requirements.txt
Zh QA (`docs/zh/QA.md`)
常见问题
为什么训练框架不支持 batch size > 1?
- 更大的 batch size 已无法实现显著加速:由于 flash attention 等加速技术已经充分提高了 GPU 的利用率,因此更大的 batch size 只会带来更大的显存占用,无法带来显著加速。在 Stable Diffusion 1.5 这类小模型上的经验已不再适用于最新的大模型。
- 更大的 batch size 可以用其他方案实现:多 GPU 训练和 Gradient Accumulation 都可以在数学意义上等价地实现更大的 batch size。
- 更大的 batch size 与框架的通用性设计相悖:我们希望构建通用的训练框架,大量模型无法适配更大的 batch size,例如不同长度的文本编码、不同分辨率的图像等,都是无法合并为更大的 batch 的。
为什么不删除某些模型中的冗余参数?
在部分模型中,模型存在冗余参数,例如 Qwen-Image 的 DiT 模型最后一层的文本部分,这部分参数不会参与任何计算,这是模型开发者留下的小 bug。直接将其设置为可训练时还会在多 GPU 训练中出现报错。
为了与开源社区中其他模型保持兼容性,我们决定保留这些参数。这些冗余参数在多 GPU 训练中可以通过 --find_unused_parameters 参数避免报错。
为什么 FP8 量化没有任何加速效果?
原生 FP8 计算需要依赖 Hopper 架构的 GPU,同时在计算精度上有较大误差,目前仍然是不成熟的技术,因此本项目不支持原生 FP8 计算。
显存管理中的 FP8 计算是指将模型参数以 FP8 精度存储在内存或显存中,在需要计算时临时转换为其他精度,因此仅能减少显存占用,没有加速效果。
为什么训练框架不支持原生 FP8 精度训练?
即使硬件条件允许,我们目前也没有任何支持原生 FP8 精度训练的规划。
- 目前原生 FP8 精度训练的主要挑战是梯度爆炸导致的精度溢出,为了保证训练的稳定性,需针对性地重新设计模型结构,然而目前还没有任何模型开发者愿意这么做。
- 此外,使用原生 FP8 精度训练的模型,在推理时若没有 Hopper 架构 GPU,则只能以 BF16 精度进行计算,理论上其生成效果反而不如 FP8。
因此,原生 FP8 精度训练技术是极不成熟的,我们静观开源社区的技术发展。
如何在推理时动态加载 LoRA 模型?
我们支持 LoRA 模型的两种加载方式,详见LoRA 加载:
- 冷加载:当基础模型未开启显存管理时,LoRA 会融合进基础模型权重,此时推理速度没有变化,LoRA 加载后无法卸载。
- 热加载:当基础模型开启显存管理时,LoRA 不会融合进基础模型权重,此时推理速度会变慢,LoRA 加载后可通过
pipe.clear_lora()卸载。
如何减少训练所需的显存?
框架支持多种方式减少训练所需的显存,包括 Gradient Checkpointing、DeepSpeed、FP8、两阶段拆分训练、CPU Offload,请参考低显存训练。
Zh Index (`docs/zh/index.rst`)
欢迎来到 DiffSynth-Studio 的文档
.. toctree::
:maxdepth: 2
:caption: 文档介绍
README
.. toctree::
:maxdepth: 2
:caption: 上手使用
Pipeline_Usage/Setup
Pipeline_Usage/Model_Inference
Pipeline_Usage/Accelerated_Inference
Pipeline_Usage/VRAM_management
Pipeline_Usage/Quantization
Pipeline_Usage/Model_Training
Pipeline_Usage/Environment_Variables
Pipeline_Usage/GPU_support
Pipeline_Usage/Inference_WebUI
.. toctree::
:maxdepth: 2
:caption: 模型详解
Model_Details/FLUX
Model_Details/Wan
Model_Details/Qwen-Image
Model_Details/Qwen-Video-Edit
Model_Details/FLUX2
Model_Details/Z-Image
Model_Details/Anima
Model_Details/LTX-2
Model_Details/LTX-2.5
Model_Details/ERNIE-Image
Model_Details/JoyAI-Image
Model_Details/ACE-Step
Model_Details/HiDream-O1-Image
Model_Details/Stable-Diffusion
Model_Details/Stable-Diffusion-XL
Model_Details/Image-Quality-Metrics
Model_Details/Ideogram-4
Model_Details/Krea-2
Model_Details/Boogu-Image
Model_Details/LingBot-Video
Model_Details/MiniMax-H3
Model_Details/DiffSynth-Music
Model_Details/MiniMax-Music3
Model_Details/YuE2
Model_Details/SenseNova-U1
Model_Details/Qwen-Image-2.1
Model_Details/FLUX.3
Model_Details/MOSS-Image
Model_Details/Fibo
.. toctree::
:maxdepth: 2
:caption: 训练框架
Training/Understanding_Diffusion_models
Training/Supervised_Fine_Tuning
Training/FP8_Precision
Training/Distillation
Training/Direct_Distill
Training/Split_Training
Training/Differential_LoRA
Training/DeepSpeed
Training/Offload_Training
.. toctree::
:maxdepth: 2
:caption: 模型接入
Developer_Guide/Integrating_Your_Model
Developer_Guide/Building_a_Pipeline
Developer_Guide/Enabling_VRAM_management
Developer_Guide/Training_Diffusion_Models
Developer_Guide/Integrating_Quantization_Backend
.. toctree::
:maxdepth: 2
:caption: API 参考
API_Reference/core/attention
API_Reference/core/data
API_Reference/core/gradient
API_Reference/core/loader
API_Reference/core/quant
API_Reference/core/vram
.. toctree::
:maxdepth: 2
:caption: Diffusion Templates
Diffusion_Templates/Introducing_Diffusion_Templates.md
Diffusion_Templates/Understanding_Diffusion_Templates.md
Diffusion_Templates/Template_Model_Inference.md
Diffusion_Templates/Template_Model_Training.md
.. toctree::
:maxdepth: 2
:caption: 学术导引
Research_Tutorial/train_from_scratch
Research_Tutorial/inference_time_scaling
Research_Tutorial/controllable_models
.. toctree::
:maxdepth: 2
:caption: 常见问题
QA
Indices and tables
- :ref:
genindex - :ref:
modindex - :ref:
search
En Developer Guide Building A Pipeline (`docs/en/Developer_Guide/Building_a_Pipeline.md`)
Building a Pipeline
After integrating the required models for the Pipeline, you also need to build a Pipeline for model inference. This document provides a standardized process for building a Pipeline. Developers can also refer to existing Pipeline implementations for construction.
The Pipeline implementation is located in diffsynth/pipelines. Each Pipeline contains the following essential key components:
__init__from_pretrained__call__unitsmodel_fn
`__init__`
In __init__, the Pipeline is initialized. Here is a simple implementation:
import torch
from PIL import Image
from typing import Union
from tqdm import tqdm
from ..diffusion import FlowMatchScheduler
from ..core import ModelConfig
from ..diffusion.base_pipeline import BasePipeline, PipelineUnit
from ..models.new_models import XXX_Model, YYY_Model, ZZZ_Model
class NewDiffSynthPipeline(BasePipeline):
def __init__(self, device="cuda", torch_dtype=torch.bfloat16):
super().__init__(device=device, torch_dtype=torch_dtype)
self.scheduler = FlowMatchScheduler()
self.text_encoder: XXX_Model = None
self.dit: YYY_Model = None
self.vae: ZZZ_Model = None
self.in_iteration_models = ("dit",)
self.units = [
NewDiffSynthPipelineUnit_xxx(),
...
]
self.model_fn = model_fn_newThis includes the following parts:
scheduler: Scheduler, used to control the coefficients in the iterative formula during inference, controlling the noise content at each step.text_encoder,dit,vae: Models. Since Latent Diffusion was proposed, this three-stage model architecture has become the mainstream Diffusion model architecture. However, this is not immutable, and any number of models can be added to thePipeline.in_iteration_models: Iteration models. This tuple marks which models will be called during iteration.units: Pre-processing units for model iteration. Seeunitsfor details.model_fn: Theforwardfunction of the denoising model during iteration. Seemodel_fnfor details.
Q: Model loading does not occur in
__init__, why initialize each model asNonehere?A: By annotating the type of each model here, the code editor can provide code completion prompts based on each model, facilitating subsequent development.
`from_pretrained`
from_pretrained is responsible for loading the required models to make the Pipeline callable. Here is a simple implementation:
@staticmethod
def from_pretrained(
torch_dtype: torch.dtype = torch.bfloat16,
device: Union[str, torch.device] = "cuda",
model_configs: list[ModelConfig] = [],
vram_limit: float = None,
):
# Initialize pipeline
pipe = NewDiffSynthPipeline(device=device, torch_dtype=torch_dtype)
model_pool = pipe.download_and_load_models(model_configs, vram_limit)
# Fetch models
pipe.text_encoder = model_pool.fetch_model("xxx_text_encoder")
pipe.dit = model_pool.fetch_model("yyy_dit")
pipe.vae = model_pool.fetch_model("zzz_vae")
# If necessary, load tokenizers here.
# VRAM Management
pipe.vram_management_enabled = pipe.check_vram_management_state()
return pipeDevelopers need to implement the logic for fetching models. The corresponding model names are the "model_name" in the model Config filled in during model integration.
Some models also need to load tokenizer. Extra tokenizer_config parameters can be added to from_pretrained as needed, and this part can be implemented after fetching the models.
`__call__`
__call__ implements the entire generation process of the Pipeline. Below is a common generation process template. Developers can modify it based on their needs.
/* Detailed source-code truncated for AI context efficiency. */`units`
units contains all the preprocessing processes, such as: width/height checking, prompt encoding, initial noise generation, etc. In the entire model preprocessing process, data is abstracted into three mutually exclusive parts, stored in corresponding dictionaries:
inputs_shared: Shared inputs, parameters unrelated to Classifier-Free Guidance (CFG for short).inputs_posi: Positive side inputs for Classifier-Free Guidance, containing content related to positive prompts.inputs_nega: Negative side inputs for Classifier-Free Guidance, containing content related to negative prompts.
Pipeline Unit implementations include three types: direct mode, CFG separation mode, and takeover mode.
If some calculations are unrelated to CFG, direct mode can be used, for example, Qwen-Image's random noise initialization:
class QwenImageUnit_NoiseInitializer(PipelineUnit):
def __init__(self):
super().__init__(
input_params=("height", "width", "seed", "rand_device"),
output_params=("noise",),
)
def process(self, pipe: QwenImagePipeline, height, width, seed, rand_device):
noise = pipe.generate_noise((1, 16, height//8, width//8), seed=seed, rand_device=rand_device, rand_torch_dtype=pipe.torch_dtype)
return {"noise": noise}If some calculations are related to CFG and need to separately process positive and negative prompts, but the input parameters on both sides are the same, CFG separation mode can be used, for example, Qwen-image's prompt encoding:
class QwenImageUnit_PromptEmbedder(PipelineUnit):
def __init__(self):
super().__init__(
seperate_cfg=True,
input_params_posi={"prompt": "prompt"},
input_params_nega={"prompt": "negative_prompt"},
input_params=("edit_image",),
output_params=("prompt_emb", "prompt_emb_mask"),
onload_model_names=("text_encoder",)
)
def process(self, pipe: QwenImagePipeline, prompt, edit_image=None) -> dict:
pipe.load_models_to_device(self.onload_model_names)
# Do something
return {"prompt_emb": prompt_embeds, "prompt_emb_mask": encoder_attention_mask}If some calculations need global information, takeover mode is required, for example, Qwen-Image's entity partition control:
class QwenImageUnit_EntityControl(PipelineUnit):
def __init__(self):
super().__init__(
take_over=True,
input_params=("eligen_entity_prompts", "width", "height", "eligen_enable_on_negative", "cfg_scale"),
output_params=("entity_prompt_emb", "entity_masks", "entity_prompt_emb_mask"),
onload_model_names=("text_encoder",)
)
def process(self, pipe: QwenImagePipeline, inputs_shared, inputs_posi, inputs_nega):
# Do something
return inputs_shared, inputs_posi, inputs_negaThe following are the parameter configurations required for Pipeline Unit:
seperate_cfg: Whether to enable CFG separation modetake_over: Whether to enable takeover modeinput_params: Shared input parametersoutput_params: Output parametersinput_params_posi: Positive side input parametersinput_params_nega: Negative side input parametersonload_model_names: Names of model components to be called
When designing unit, please try to follow these principles:
- Default fallback: For optional function
unitinput parameters, the default isNonerather thanFalseor other values. Please provide fallback processing for this default value. - Parameter triggering: Some Adapter models may not be loaded, such as ControlNet. The corresponding
unitshould control triggering based on whether the parameter input isNonerather than whether the model is loaded. For example, when the user inputscontrolnet_imagebut does not load the ControlNet model, the code should give an error rather than ignore these input parameters and continue execution. - Simplicity first: Use direct mode as much as possible, only use takeover mode when the function cannot be implemented.
- VRAM efficiency: When calling models in
unit, please usepipe.load_models_to_device(self.onload_model_names)to activate the corresponding models. Do not call other models outsideonload_model_names. Afterunitcalculation is completed, do not manually release VRAM withpipe.load_models_to_device([]).
Q: Some parameters are not called during the inference process, such as
output_params. Is it still necessary to configure them?A: These parameters will not affect the inference process, but they will affect some experimental features. Therefore, we recommend configuring them properly. For example, "split training" - we can complete the preprocessing offline during training, but some model calculations that require gradient backpropagation cannot be split. These parameters are used to build computational graphs to infer which calculations can be split.
`model_fn`
model_fn is the unified forward interface during iteration. For models where the open-source ecosystem is not yet formed, you can directly use the denoising model's forward, for example:
def model_fn_new(dit=None, latents=None, timestep=None, prompt_emb=None, **kwargs):
return dit(latents, prompt_emb, timestep)For models with rich open-source ecosystems, model_fn usually contains complex and chaotic cross-model inference. Taking diffsynth/pipelines/qwen_image.py as an example, the additional calculations implemented in this function include: entity partition control, three types of ControlNet, Gradient Checkpointing, etc. Developers need to be extra careful when implementing this part to avoid conflicts between module functions.
Compilation Acceleration
To enable compilation acceleration, please refer to Inference Acceleration.
En Developer Guide Enabling VRAM Management (`docs/en/Developer_Guide/Enabling_VRAM_management.md`)
Fine-Grained VRAM Management Scheme
This document introduces how to write reasonable fine-grained VRAM management schemes for models, and how to use the VRAM management functions in DiffSynth-Studio for other external code libraries. Before reading this document, please read the document VRAM Management.
How Much VRAM Does a 20B Model Need?
Taking Qwen-Image's DiT model as an example, this model has reached 20B parameters. The following code will load this model and perform inference, requiring about 40G VRAM. This model obviously cannot run on consumer-grade GPUs with smaller VRAM.
from diffsynth.core import load_model
from diffsynth.models.qwen_image_dit import QwenImageDiT
from modelscope import snapshot_download
import torch
snapshot_download(
model_id="Qwen/Qwen-Image",
local_dir="models/Qwen/Qwen-Image",
allow_file_pattern="transformer/*"
)
prefix = "models/Qwen/Qwen-Image/transformer/diffusion_pytorch_model"
model_path = [prefix + f"-0000{i}-of-00009.safetensors" for i in range(1, 10)]
inputs = {
"latents": torch.randn((1, 16, 128, 128), dtype=torch.bfloat16, device="cuda"),
"timestep": torch.zeros((1,), dtype=torch.bfloat16, device="cuda"),
"prompt_emb": torch.randn((1, 5, 3584), dtype=torch.bfloat16, device="cuda"),
"prompt_emb_mask": torch.ones((1, 5), dtype=torch.int64, device="cuda"),
"height": 1024,
"width": 1024,
}
model = load_model(QwenImageDiT, model_path, torch_dtype=torch.bfloat16, device="cuda")
with torch.no_grad():
output = model(**inputs)Writing Fine-Grained VRAM Management Scheme
To write a fine-grained VRAM management scheme, we need to use print(model) to observe and analyze the model structure:
/* Detailed source-code truncated for AI context efficiency. */In VRAM management, we only care about layers containing parameters. In this model structure, QwenEmbedRope, TemporalTimesteps, SiLU and other Layers do not contain parameters. LayerNorm also does not contain parameters because elementwise_affine=False is set. Layers containing parameters are only Linear and RMSNorm.
diffsynth.core.vram provides two replacement modules for VRAM management:
AutoWrappedLinear: Used to replaceLinearlayersAutoWrappedModule: Used to replace any other layer
Write a module_map to map Linear and RMSNorm in the model to the corresponding modules:
module_map={
torch.nn.Linear: AutoWrappedLinear,
RMSNorm: AutoWrappedModule,
}In addition, vram_config and vram_limit are also required, which have been introduced in VRAM Management.
Call enable_vram_management to enable VRAM management. Note that the device when loading the model is cpu, consistent with offload_device:
from diffsynth.core import load_model, enable_vram_management, AutoWrappedLinear, AutoWrappedModule
from diffsynth.models.qwen_image_dit import QwenImageDiT, RMSNorm
import torch
prefix = "models/Qwen/Qwen-Image/transformer/diffusion_pytorch_model"
model_path = [prefix + f"-0000{i}-of-00009.safetensors" for i in range(1, 10)]
inputs = {
"latents": torch.randn((1, 16, 128, 128), dtype=torch.bfloat16, device="cuda"),
"timestep": torch.zeros((1,), dtype=torch.bfloat16, device="cuda"),
"prompt_emb": torch.randn((1, 5, 3584), dtype=torch.bfloat16, device="cuda"),
"prompt_emb_mask": torch.ones((1, 5), dtype=torch.int64, device="cuda"),
"height": 1024,
"width": 1024,
}
model = load_model(QwenImageDiT, model_path, torch_dtype=torch.bfloat16, device="cpu")
enable_vram_management(
model,
module_map={
torch.nn.Linear: AutoWrappedLinear,
RMSNorm: AutoWrappedModule,
},
vram_config = {
"offload_dtype": torch.bfloat16,
"offload_device": "cpu",
"onload_dtype": torch.bfloat16,
"onload_device": "cpu",
"preparing_dtype": torch.bfloat16,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
},
vram_limit=0,
)
with torch.no_grad():
output = model(**inputs)The above code only requires 2G VRAM to run the forward of a 20B model.
Disk Offload
Disk Offload is a special VRAM management scheme that needs to be enabled during the model loading process, not after the model is loaded. Usually, when the above code can run smoothly, Disk Offload can be directly enabled:
from diffsynth.core import load_model, enable_vram_management, AutoWrappedLinear, AutoWrappedModule
from diffsynth.models.qwen_image_dit import QwenImageDiT, RMSNorm
import torch
prefix = "models/Qwen/Qwen-Image/transformer/diffusion_pytorch_model"
model_path = [prefix + f"-0000{i}-of-00009.safetensors" for i in range(1, 10)]
inputs = {
"latents": torch.randn((1, 16, 128, 128), dtype=torch.bfloat16, device="cuda"),
"timestep": torch.zeros((1,), dtype=torch.bfloat16, device="cuda"),
"prompt_emb": torch.randn((1, 5, 3584), dtype=torch.bfloat16, device="cuda"),
"prompt_emb_mask": torch.ones((1, 5), dtype=torch.int64, device="cuda"),
"height": 1024,
"width": 1024,
}
model = load_model(
QwenImageDiT,
model_path,
module_map={
torch.nn.Linear: AutoWrappedLinear,
RMSNorm: AutoWrappedModule,
},
vram_config={
"offload_dtype": "disk",
"offload_device": "disk",
"onload_dtype": "disk",
"onload_device": "disk",
"preparing_dtype": torch.bfloat16,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
},
vram_limit=0,
)
with torch.no_grad():
output = model(**inputs)Disk Offload is an extremely special VRAM management scheme. It only supports .safetensors format files, not binary files such as .bin, .pth, .ckpt, and does not support state dict converter with Tensor reshape.
If there are situations where Disk Offload cannot run normally but non-Disk Offload can run normally, please submit an issue to us on GitHub.
Writing Default Configuration
To make it easier for users to use the VRAM management function, we write the fine-grained VRAM management configuration in diffsynth/configs/vram_management_module_maps.py. The configuration information for the above model is:
"diffsynth.models.qwen_image_dit.QwenImageDiT": {
"diffsynth.models.qwen_image_dit.RMSNorm": "diffsynth.core.vram.layers.AutoWrappedModule",
"torch.nn.Linear": "diffsynth.core.vram.layers.AutoWrappedLinear",
}
```# Fine-Grained VRAM Management Scheme
This document introduces how to write reasonable fine-grained VRAM management schemes for models, and how to use the VRAM management functions in `DiffSynth-Studio` for other external code libraries. Before reading this document, please read the document [VRAM Management](../Pipeline_Usage/VRAM_management.md).
## How Much VRAM Does a 20B Model Need?
Taking Qwen-Image's DiT model as an example, this model has reached 20B parameters. The following code will load this model and perform inference, requiring about 40G VRAM. This model obviously cannot run on consumer-grade GPUs with smaller VRAM.
```python
from diffsynth.core import load_model
from diffsynth.models.qwen_image_dit import QwenImageDiT
from modelscope import snapshot_download
import torch
snapshot_download(
model_id="Qwen/Qwen-Image",
local_dir="models/Qwen/Qwen-Image",
allow_file_pattern="transformer/*"
)
prefix = "models/Qwen/Qwen-Image/transformer/diffusion_pytorch_model"
model_path = [prefix + f"-0000{i}-of-00009.safetensors" for i in range(1, 10)]
inputs = {
"latents": torch.randn((1, 16, 128, 128), dtype=torch.bfloat16, device="cuda"),
"timestep": torch.zeros((1,), dtype=torch.bfloat16, device="cuda"),
"prompt_emb": torch.randn((1, 5, 3584), dtype=torch.bfloat16, device="cuda"),
"prompt_emb_mask": torch.ones((1, 5), dtype=torch.int64, device="cuda"),
"height": 1024,
"width": 1024,
}
model = load_model(QwenImageDiT, model_path, torch_dtype=torch.bfloat16, device="cuda")
with torch.no_grad():
output = model(**inputs)Writing Fine-Grained VRAM Management Scheme
To write a fine-grained VRAM management scheme, we need to use print(model) to observe and analyze the model structure:
/* Detailed source-code truncated for AI context efficiency. */In VRAM management, we only care about layers containing parameters. In this model structure, QwenEmbedRope, TemporalTimesteps, SiLU and other Layers do not contain parameters. LayerNorm also does not contain parameters because elementwise_affine=False is set. Layers containing parameters are only Linear and RMSNorm.
diffsynth.core.vram provides two replacement modules for VRAM management:
AutoWrappedLinear: Used to replaceLinearlayersAutoWrappedModule: Used to replace any other layer
Write a module_map to map Linear and RMSNorm in the model to the corresponding modules:
module_map={
torch.nn.Linear: AutoWrappedLinear,
RMSNorm: AutoWrappedModule,
}In addition, vram_config and vram_limit are also required, which have been introduced in VRAM Management.
Call enable_vram_management to enable VRAM management. Note that the device when loading the model is cpu, consistent with offload_device:
from diffsynth.core import load_model, enable_vram_management, AutoWrappedLinear, AutoWrappedModule
from diffsynth.models.qwen_image_dit import QwenImageDiT, RMSNorm
import torch
prefix = "models/Qwen/Qwen-Image/transformer/diffusion_pytorch_model"
model_path = [prefix + f"-0000{i}-of-00009.safetensors" for i in range(1, 10)]
inputs = {
"latents": torch.randn((1, 16, 128, 128), dtype=torch.bfloat16, device="cuda"),
"timestep": torch.zeros((1,), dtype=torch.bfloat16, device="cuda"),
"prompt_emb": torch.randn((1, 5, 3584), dtype=torch.bfloat16, device="cuda"),
"prompt_emb_mask": torch.ones((1, 5), dtype=torch.int64, device="cuda"),
"height": 1024,
"width": 1024,
}
model = load_model(QwenImageDiT, model_path, torch_dtype=torch.bfloat16, device="cpu")
enable_vram_management(
model,
module_map={
torch.nn.Linear: AutoWrappedLinear,
RMSNorm: AutoWrappedModule,
},
vram_config = {
"offload_dtype": torch.bfloat16,
"offload_device": "cpu",
"onload_dtype": torch.bfloat16,
"onload_device": "cpu",
"preparing_dtype": torch.bfloat16,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
},
vram_limit=0,
)
with torch.no_grad():
output = model(**inputs)The above code only requires 2G VRAM to run the forward of a 20B model.
Disk Offload
Disk Offload is a special VRAM management scheme that needs to be enabled during the model loading process, not after the model is loaded. Usually, when the above code can run smoothly, Disk Offload can be directly enabled:
from diffsynth.core import load_model, enable_vram_management, AutoWrappedLinear, AutoWrappedModule
from diffsynth.models.qwen_image_dit import QwenImageDiT, RMSNorm
import torch
prefix = "models/Qwen/Qwen-Image/transformer/diffusion_pytorch_model"
model_path = [prefix + f"-0000{i}-of-00009.safetensors" for i in range(1, 10)]
inputs = {
"latents": torch.randn((1, 16, 128, 128), dtype=torch.bfloat16, device="cuda"),
"timestep": torch.zeros((1,), dtype=torch.bfloat16, device="cuda"),
"prompt_emb": torch.randn((1, 5, 3584), dtype=torch.bfloat16, device="cuda"),
"prompt_emb_mask": torch.ones((1, 5), dtype=torch.int64, device="cuda"),
"height": 1024,
"width": 1024,
}
model = load_model(
QwenImageDiT,
model_path,
module_map={
torch.nn.Linear: AutoWrappedLinear,
RMSNorm: AutoWrappedModule,
},
vram_config={
"offload_dtype": "disk",
"offload_device": "disk",
"onload_dtype": "disk",
"onload_device": "disk",
"preparing_dtype": torch.bfloat16,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
},
vram_limit=0,
)
with torch.no_grad():
output = model(**inputs)Disk Offload is an extremely special VRAM management scheme. It only supports .safetensors format files, not binary files such as .bin, .pth, .ckpt, and does not support state dict converter with Tensor reshape.
If there are situations where Disk Offload cannot run normally but non-Disk Offload can run normally, please submit an issue to us on GitHub.
Writing Default Configuration
To make it easier for users to use the VRAM management function, we write the fine-grained VRAM management configuration in diffsynth/configs/vram_management_module_maps.py. The configuration information for the above model is:
"diffsynth.models.qwen_image_dit.QwenImageDiT": {
"diffsynth.models.qwen_image_dit.RMSNorm": "diffsynth.core.vram.layers.AutoWrappedModule",
"torch.nn.Linear": "diffsynth.core.vram.layers.AutoWrappedLinear",
}En Developer Guide Integrating Quantization Backend (`docs/en/Developer_Guide/Integrating_Quantization_Backend.md`)
Integrating a Quantization Backend
The quantization framework of DiffSynth-Studio lives in diffsynth.core.quant and ships with bitsandbytes, torchao, and comfy-kitchen backends (see Model Quantization). If you have your own quantization algorithm, or want to plug in another quantization library, you only need to implement a QuantBackend — online quantization, saving/loading pre-quantized checkpoints, mixed quantization, VRAM management, and quantization + LoRA training are all reused as-is.
This guide walks through the whole process with a toy backend: INT9 — 9-bit symmetric weight-only quantization, genuinely stored at 9 bits per weight, with one fp32 scale per output channel. INT9 does not exist on any hardware; it is used here because it keeps the example short while still covering every interface you have to implement. For the full interface signatures and contracts, see the diffsynth.core.quant API documentation.
Framework Structure
The framework has three layers:
QuantizeConfig: the user-facing config and entry point, responsible for traversing the model, matching layers, and replacingnn.Linear. You never need to touch it.QuantBackend: the adapter layer, which only ever deals with a singlenn.Linear: how to quantize it, how to build an empty shell, how to dequantize it, how to read and write its state dict. This is the part you implement.- The quantized Linear: the module that actually holds the quantized weight and performs dequantization + matmul in
forward.
The quantized Linear must satisfy four contract clauses:
- (a) It is a drop-in replacement for
nn.Linear, withforward(x)doing dequantization + matmul internally. It must subclasstorch.nn.Linear, otherwise LoRA injection and VRAM management cannot see it. - (b)
.to(...)only moves devices, never re-types the packed weight or quant state. VRAM management performs dtype conversions on the model; if a packed weight were cast to bf16, the quant state would be corrupted. - (c)
state_dict()andload_state_dict(assign=True)round-trip, viaflatten_state_dict/unflatten_state_dictwhen necessary. - (d) (Training only)
forwardis differentiable w.r.t. its input, so gradients can pass through frozen quantized layers to reach LoRA branches.
Step 1: Write the Quantized Linear
INT9's storage layout needs a little thought: there is no native 9-bit dtype, and simply putting the codes into an int16 tensor would still spend 16 bits per weight — exactly as much as bf16, so the quantization would save nothing. Each weight is therefore split in two: the low 8 bits go into a uint8 weight buffer, and the 9th (most significant) bit forms a separate bit plane where 8 weights are packed into one byte in weight_msb, plus one fp32 weight_scale per output channel. That is 9 bits per weight, 56% of bf16.
Two more details matter:
- Delete
nn.Linear's originalweightparameter and register a buffer under the same name, so checkpoint keys staylayer_name.weight. Disk offload and mixed-quantization key ownership rely on this. - Guard the dtype of the packed tensors by overriding
_apply, i.e. contract clause (b). Every conversion (.to(),.half(),.float(), ...) funnels through_apply, so a conversion that would change the dtype is downgraded to a device-only move.
/* Detailed source-code truncated for AI context efficiency. */Dequantization in forward uses ordinary tensor ops, so gradients flow back to the input x through F.linear and contract clause (d) holds automatically. Unpacking here is written bit by bit with PyTorch ops purely for clarity; a real backend fuses unpacking into the matmul kernel instead of materializing an fp weight on every forward.
There is an easy trap in dequantize_weight: the integer codes must be reconstructed in fp32. bf16 only carries 8 bits of significand, so integers above 256 are not representable exactly; casting the codes to bf16 before applying the scale rounds the 9th bit away and throws the accuracy gain out (measured: the advantage over int8 collapses from 2.25x to 1.15x). This applies to any format whose code width exceeds the significand of the compute dtype.
Step 2: Write the Backend
Every backend method operates on a single nn.Linear:
capabilities(): declares what the backend supports; all four flags default toFalse. Saving quantized weights requiresis_serializable=True, and quantization + LoRA training requiresis_differentiable=True.quantized_linear_classes(): declares the Linear classes this backend produces;is_quantized_lineardefaults to anisinstancecheck against them.create_quantized_linear(): online quantization, turning an fpnn.Linearinto a quantized one.compute_deviceis where quantization runs andmodel_deviceis where the result is stored, so the two together stream the work layer by layer with only one layer on the accelerator at a time.create_quantized_linear_shell(): builds an empty shell, used for loading pre-quantized checkpoints and for disk offload. It is rebuilt on every offload cycle, so build it on themetadevice and keep it cheap.dequantize_to_linear(): restores a plainnn.Linear, used bymode="dequant_once".flatten_state_dict/unflatten_state_dict: conversion between the state dict and flat tensors. INT9's state dict already holds plain tensors, so the base class implementation is enough; only backends with composite tensors (tensor subclasses, nested quant state) such as bitsandbytes and torchao need to override them.
Unimplemented methods raise a descriptive exception from the base class, so a backend that only supports some capabilities can implement just what it needs. self.config is the backend config instance injected by the framework — the Int9WeightOnlyConfig written in the next step.
/* Detailed source-code truncated for AI context efficiency. */Step 3: Write the Backend Config
The backend config subclasses BackendConfig: user-tunable parameters are ordinary dataclass fields, while values pinned by the method are declared with field(init=False, default=...). describe_quant_method reports the two groups separately, and from_kwargs raises when a user passes unknown backend_config_kwargs.
@dataclass
class Int9WeightOnlyConfig(BackendConfig):
per_channel: bool = True # user-tunable: per-channel or per-tensor
bits: int = field(init=False, default=9) # pinned by the method, not overridableStep 4: Register the Quantization Method
One backend can register several methods, distinguished by the fields pinned in its config (the bitsandbytes backend, for example, distinguishes nf4 from fp4 via quant_type). Method names should follow the _<format>_w<weight bits>a<activation bits> convention:
register_quant_method("toy_int9_w9a16", "toy_int9", Int9WeightOnlyConfig.from_kwargs, label="9bit, int9, weight-only (toy)")There are two ways to register a backend and its methods:
Option 1: keep it in your own code (recommended, plug-and-play). Put the code above in any module; as long as that module is imported before you construct QuantizeConfig, the method is already in QUANT_METHODS and can be used just like a built-in one, with no framework changes:
import my_project.toy_int9 # triggers register_quant_backend / register_quant_method
from diffsynth.core.quant import QuantizeConfig
quantize = QuantizeConfig(method="toy_int9_w9a16", backend_config_kwargs={"per_channel": True})Option 2: ship it as a built-in backend (permanent). Put the backend file under diffsynth/core/quant/backends/ and register it in _LAZY_BACKENDS in diffsynth/core/quant/backends/__init__.py; the framework then imports it on demand and users do not need to import anything:
_LAZY_BACKENDS = {
"bitsandbytes": ".bitsandbytes",
"torchao": ".torchao",
"comfy_kitchen": ".comfy_kitchen",
"toy_int9": ".toy_int9",
}If your quantization algorithm or library is generally useful, you are welcome to submit it as a PR following Option 2, so that more users can benefit from it. A backend that depends on a third-party library should check its dependencies in validate_environment() with an installation hint, and point project_url at the upstream project.
Step 5: Self-Check
The framework provides two verification tools; run them right after integrating. check_backend_contract verifies that the backend declares its Linear classes, that both factory methods return instances of those classes, that all declared classes subclass nn.Linear, and that every checkpoint key the backend actually writes lives under the layer name (a missing scale would make disk offload silently load corrupted layers). Unsupported factory methods are skipped rather than counted as failures.
from diffsynth.core.quant import QUANT_BACKENDS, QUANT_METHODS, check_backend_contract, check_differentiable, describe_quant_method
describe_quant_method("toy_int9_w9a16")
spec = QUANT_METHODS["toy_int9_w9a16"]
check_backend_contract(QUANT_BACKENDS[spec.backend](spec.config_factory({})), compute_device="cpu")The output is as follows; describe_quant_method also confirms that the split between user-tunable and pinned parameters is what you intended:
method: toy_int9_w9a16
backend: toy_int9
detail: 9bit, int9, weight-only (toy)
backend config: my_project.toy_int9.Int9WeightOnlyConfig
backend_config_kwargs (user-tunable):
per_channel = True
pinned by method (not overridable):
bits = 9
check_backend_contract (toy_int9):
[PASS] quantized_linear_classes() is non-empty: ['Int9Linear']
[PASS] Int9Linear subclasses torch.nn.Linear
[PASS] a plain nn.Linear is not reported as quantized
[PASS] create_quantized_linear_shell() returns a declared class, got Int9Linear
[PASS] the shell is recognized before load_state_dict (disk offload routing)
[PASS] create_quantized_linear() returns a declared class, got Int9Linear
[PASS] every stored key lives under the layer name; uncovered: []
=> OKNext, check the numerical error, the real memory saving, the dtype guard of clause (b), and the differentiability of clause (d) on a small model:
import torch
from diffsynth.core.quant import QuantizeConfig, check_differentiable
class ToyModel(torch.nn.Module):
def __init__(self):
super().__init__()
self.fc1 = torch.nn.Linear(256, 512)
self.fc2 = torch.nn.Linear(512, 256, bias=False)
def forward(self, x):
return self.fc2(torch.nn.functional.silu(self.fc1(x)))
def footprint(model):
return sum(t.numel() * t.element_size() for t in list(model.parameters()) + list(model.buffers()))
torch.manual_seed(0)
model = ToyModel().to(torch.bfloat16)
x = torch.randn(4, 256, dtype=torch.bfloat16)
reference = model(x)
fp_bytes = footprint(model)
QuantizeConfig(method="toy_int9_w9a16").quantize_model(model, compute_device="cpu")
print("relative error:", ((model(x) - reference).norm() / reference.norm()).item())
print(f"footprint: {fp_bytes} -> {footprint(model)} bytes ({footprint(model) / fp_bytes:.3f} of bf16)")
model.to(torch.float32) # clause (b): packed dtypes must not change
print(model.fc1.weight.dtype, model.fc1.weight_msb.dtype, model.fc1.weight_scale.dtype, model.fc1.bias.dtype)
check_differentiable(model.fc1) # clause (d)2 nn.Linear layers quantized (method: toy_int9_w9a16).
relative error: 0.004150390625
footprint: 525312 -> 299008 bytes (0.569 of bf16)
torch.uint8 torch.uint8 torch.float32 torch.float32
check_differentiable (Int9Linear): OK -- gradients pass through the module to its inputThe measured footprint is 0.569 of bf16, slightly above 9/16 = 0.5625 because of the fp32 scales and the unquantized bias. If this ratio comes out close to 1, the packing format is not actually compressing the weights and you should revisit the storage layout in Step 1.
Inference on a Real Model: Z-Image
Once the small-model checks pass, the backend is ready for real models — a custom backend is used exactly like a built-in method. Import the module that registers it, then pass the method to ModelConfig(quantize=...):
import torch
import my_project.toy_int9 # registers the toy_int9 backend and the toy_int9_w9a16 method
from diffsynth.core.quant import QuantizeConfig
from diffsynth.pipelines.z_image import ModelConfig, ZImagePipeline
pipe = ZImagePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(
model_id="Tongyi-MAI/Z-Image-Turbo",
origin_file_pattern="transformer/*.safetensors",
quantize=QuantizeConfig(method="toy_int9_w9a16"),
),
ModelConfig(model_id="Tongyi-MAI/Z-Image-Turbo", origin_file_pattern="text_encoder/*.safetensors"),
ModelConfig(model_id="Tongyi-MAI/Z-Image-Turbo", origin_file_pattern="vae/diffusion_pytorch_model.safetensors"),
],
tokenizer_config=ModelConfig(model_id="Tongyi-MAI/Z-Image-Turbo", origin_file_pattern="tokenizer/"),
)
dit_bytes = sum(t.numel() * t.element_size() for t in list(pipe.dit.parameters()) + list(pipe.dit.buffers()))
print(f"dit weights: {dit_bytes / 1024 ** 3:.3f} GiB")
prompt = "A delicate portrait of an underwater girl, blue dress flowing, hair gently drifting, light and shadow clear, surrounded by bubbles, serene expression, exquisite details, dreamlike and beautiful."
image = pipe(prompt=prompt, seed=42, rand_device="cuda")
image.save("z_image_toy_int9.jpg")Measured DiT weight footprint on Z-Image Turbo (the 8-step Turbo generation works normally, with no visible quality difference from bf16):
| DiT weights | |
|---|---|
| bf16 | 11.464 GiB |
toy_int9_w9a16 |
6.456 GiB (0.563x) |
Note that peak memory and weight footprint are not the same thing: this toy materializes a temporary fp weight on every forward, so the peak saving is smaller than the storage saving. Measured on a synthetic model with 48 Linears, all weights resident on the GPU:
| Weights | Forward peak | |
|---|---|---|
| bf16 | 1.500 GiB | 1.527 GiB |
toy_int9_w9a16 |
0.845 GiB (0.563x) | 1.036 GiB (0.678x) |
That temporary weight depends only on the largest single layer and does not grow with depth, so the deeper the model, the closer the peak saving gets to the weight ratio; a real backend fusing unpacking into the matmul kernel does not need it at all. To push the peak down further, stack VRAM management on top and move weights layer by layer (pass a vram_config to each ModelConfig above — measured peak drops to 2.1 GiB).
Accuracy: int9 vs int8
Does the extra bit actually buy accuracy? Quantize the same weight with the identical per-channel symmetric scheme at 8 and 9 bits, then compare the dequantized weight error and the layer output error. This is the general recipe for an accuracy regression on a new backend: hold everything else fixed and change only the bit width.
import torch
from my_project.toy_int9 import Int9QuantBackend, Int9WeightOnlyConfig
def quantize_int8(linear):
"""The same per-channel symmetric scheme with one bit less: codes in [-128, 127]."""
weight = linear.weight.data
scale = (weight.abs().amax(dim=1).float() / 127).clamp(min=1e-8)
codes = (weight.float() / scale.unsqueeze(1)).round().clamp(-128, 127)
return (codes * scale.unsqueeze(1)).to(weight.dtype)
def relative_error(reference, value):
return ((value.float() - reference.float()).norm() / reference.float().norm()).item()
torch.manual_seed(0)
backend = Int9QuantBackend(Int9WeightOnlyConfig())
linear = torch.nn.Linear(2048, 2048, bias=False).to(torch.bfloat16)
fp_weight = linear.weight.data.clone()
int9_weight = backend.create_quantized_linear(linear).dequantize_weight(torch.bfloat16)
int8_weight = quantize_int8(linear)
error8, error9 = relative_error(fp_weight, int8_weight), relative_error(fp_weight, int9_weight)
print(f"weight error: int8 {error8:.6f} | int9 {error9:.6f} ({error8 / error9:.2f}x lower)")
x = torch.randn(64, 2048, dtype=torch.bfloat16)
reference = torch.nn.functional.linear(x, fp_weight)
out8 = relative_error(reference, torch.nn.functional.linear(x, int8_weight))
out9 = relative_error(reference, torch.nn.functional.linear(x, int9_weight))
print(f"output error: int8 {out8:.6f} | int9 {out9:.6f} ({out8 / out9:.2f}x lower)")weight error: int8 0.004353 | int9 0.001937 (2.25x lower)
output error: int8 0.004947 | int9 0.002816 (1.76x lower)This matches the theory: going from 255 to 511 levels halves the quantization step, and for uniform quantization the error is proportional to the step, so the weight error drops to roughly half (2.25x measured). The end-to-end layer output gain is smaller (1.76x) because the activations themselves are bf16 and the matmul's own rounding noise eats part of the benefit — a reminder to evaluate bit-width gains at the actual compute precision, not only on the weights.
Finally, verify clause (c): save the quantized weights, load them back into shells, and confirm both produce identical outputs.
from safetensors.torch import load_file, save_file
save_config = QuantizeConfig(method="toy_int9_w9a16")
tensors, metadata = save_config.flatten_state_dict(model.state_dict())
save_file(tensors, "toy_int9.safetensors", metadata=metadata)
loaded = ToyModel().to(torch.bfloat16)
load_config = QuantizeConfig(method="toy_int9_w9a16", load_prequantized=True)
load_config.prepare_for_prequantized_load(loaded, compute_dtype=torch.bfloat16)
loaded.load_state_dict(load_config.unflatten_state_dict(load_file("toy_int9.safetensors"), metadata), assign=True)
print("reload match:", torch.equal(loaded(x.float()), model(x.float())))reload match: TrueCombining with Disk Offload
Disk offload, part of VRAM management, places the strictest demands on a quantization backend: the resident model keeps only meta shells, and each layer's tensors are streamed back from disk at forward time and dropped right after. It relies on two things:
- It only supports pre-quantized checkpoints, so
load_prequantized=Trueis required andprepare_for_prequantized_loadmust first swap the target layers for shells. - Which tensors a layer needs is resolved by a prefix scan over the checkpoint keys using the layer's dotted name, and the result is loaded with a strict
load_state_dict(assign=True). So the only requirement on a backend is that every tensor lives under{layer_name}.— flat siblings likelayer.weight_scaleand nested quant state like bnb's both work. A missing or extra key raises instead of silently loading a corrupted layer.
/* Detailed source-code truncated for AI context efficiency. */Measured on the same ToyModel (torch_dtype=None on DiskMap is essential — it guarantees the packed tensors are not re-typed while being read):
2 nn.Linear layers replaced for loading the pre-quantized checkpoint (method: toy_int9_w9a16).
fc1: ['fc1.bias', 'fc1.weight', 'fc1.weight_msb', 'fc1.weight_scale']
fc2: ['fc2.weight', 'fc2.weight_msb', 'fc2.weight_scale']
resident 299008 bytes -> in memory after disk offload 0 bytes
output matches: True | repeatable: TrueEach layer's weight / weight_msb / weight_scale / bias is correctly attributed to that layer, the resident footprint drops to 0 bytes (everything is a meta shell), the output is bit-identical to the resident quantized model, and repeated forwards stay stable — so rebuilding shells and streaming from disk has no side effects.
On a real model, use the standard workflows from Model Quantization for end-to-end validation: pass QuantizeConfig(method="toy_int9_w9a16") to ModelConfig(quantize=...) for online-quantized inference, save the quantized weights with save_quantized_model and load them back after registering the hash, and inject LoRA into the quantized model for training.
Integration Checklist
- The packing format really shrinks the weights: the measured footprint ratio should be close to the theoretical bit-width ratio, not close to 1.
- The accuracy gain is verified: compared against the same scheme with one bit less, the error really goes down; otherwise precision is being lost somewhere in the dequantization path.
- The quantized Linear subclasses
torch.nn.Linear, and allstate_dictkeys live under the layer name. _applyguards the dtype of every packed tensor and quant state.capabilities()matches reality: declaringis_serializablerequires a round-tripping state dict, and declaringis_differentiablerequires passingcheck_differentiable.create_quantized_linearhonorscompute_device/model_device, so layer-by-layer streaming quantization works.- Disk offload works: the shell is built on
metaand cheap to rebuild, every stored tensor lives under the layer's dotted name, andunflatten_state_dicttolerates being called with a single layer's subdict plus whole-file metadata. - When depending on a third-party library,
validate_environment()gives a clear installation hint andproject_urlpoints at the upstream project. check_backend_contractpasses completely.
En Developer Guide Integrating Your Model (`docs/en/Developer_Guide/Integrating_Your_Model.md`)
Integrating Model Architecture
This document introduces how to integrate models into the DiffSynth-Studio framework for use by modules such as Pipeline.
Step 1: Integrate Model Architecture Code
All model architecture implementations in DiffSynth-Studio are unified in diffsynth/models. Each .py code file implements a model architecture, and all models are loaded through ModelPool in diffsynth/models/model_loader.py. When integrating new model architectures, please create a new .py file under this path.
diffsynth/models/
├── general_modules.py
├── model_loader.py
├── qwen_image_controlnet.py
├── qwen_image_dit.py
├── qwen_image_text_encoder.py
├── qwen_image_vae.py
└── ...In most cases, we recommend integrating models in native PyTorch code form, with the model architecture class directly inheriting from torch.nn.Module, for example:
import torch
class NewDiffSynthModel(torch.nn.Module):
def __init__(self, dim=1024):
super().__init__()
self.linear = torch.nn.Linear(dim, dim)
self.activation = torch.nn.Sigmoid()
def forward(self, x):
x = self.linear(x)
x = self.activation(x)
return xIf the model architecture implementation contains additional dependencies, we strongly recommend removing them, otherwise this will cause heavy package dependency issues. In our existing models, Qwen-Image's Blockwise ControlNet is integrated in this way. The code is lightweight, please refer to diffsynth/models/qwen_image_controlnet.py.
If the model has been integrated by Huggingface Library (transformers, diffusers, etc.), we can integrate the model in a simpler way:
Integrating Huggingface Library Style Model Architecture Code
The loading method for these models in Huggingface Library is:
from transformers import XXX_Model
model = XXX_Model.from_pretrained("path_to_your_model")DiffSynth-Studio does not support loading models through from_pretrained because this conflicts with VRAM management and other functions. Please rewrite the model architecture in the following format:
import torch
class DiffSynth_XXX_Model(torch.nn.Module):
def __init__(self):
super().__init__()
from transformers import XXX_Config, XXX_Model
config = XXX_Config(**{
"architectures": ["XXX_Model"],
"other_configs": "Please copy and paste the other configs here.",
})
self.model = XXX_Model(config)
def forward(self, x):
outputs = self.model(x)
return outputsWhere XXX_Config is the Config class corresponding to the model. For example, the Config class for Qwen2_5_VLModel is Qwen2_5_VLConfig, which can be found by consulting its source code. The content inside Config can usually be found in the config.json file in the model library. DiffSynth-Studio will not read the config.json file, so the content needs to be copied and pasted into the code.
In rare cases, version updates of transformers and diffusers may cause some models to be unable to import. Therefore, if possible, we still recommend using the model integration method in Step 1.1.
In our existing models, Qwen-Image's Text Encoder is integrated in this way. The code is lightweight, please refer to diffsynth/models/qwen_image_text_encoder.py.
Step 2: Model File Format Conversion
Due to the variety of model file formats provided by developers in the open-source community, we sometimes need to convert model file formats to form correctly formatted state dict. This is common in the following situations:
- Model files built by different code libraries, for example Wan-AI/Wan2.1-T2V-1.3B and Wan-AI/Wan2.1-T2V-1.3B-Diffusers.
- Models modified during integration, for example, the Text Encoder of Qwen/Qwen-Image adds a
model.prefix indiffsynth/models/qwen_image_text_encoder.py. - Model files containing multiple models, for example, the VACE Adapter and base DiT model of Wan-AI/Wan2.1-VACE-14B are mixed and stored in the same set of model files.
In our development philosophy, we hope to respect the wishes of model authors as much as possible. If we repackage the model files, for example Comfy-Org/Qwen-Image_ComfyUI, although we can call the model more conveniently, traffic (model page views and downloads, etc.) will be directed elsewhere, and the original author of the model will also lose the power to delete the model. Therefore, we have added the diffsynth/utils/state_dict_converters module to the framework for file format conversion during model loading.
This part of logic is very simple. Taking Qwen-Image's Text Encoder as an example, only 10 lines of code are needed:
def QwenImageTextEncoderStateDictConverter(state_dict):
state_dict_ = {}
for k in state_dict:
v = state_dict[k]
if k.startswith("visual."):
k = "model." + k
elif k.startswith("model."):
k = k.replace("model.", "model.language_model.")
state_dict_[k] = v
return state_dict_Step 3: Writing Model Config
Model Config is located in diffsynth/configs/model_configs.py, used to identify model types and load them. The following fields need to be filled in:
model_hash: Model file hash value, which can be obtained through thehash_model_filefunction. This hash value is only related to the keys and tensor shapes in the model file's state dict, and is unrelated to other information in the file.model_name: Model name, used forPipelineto identify the required model. If different structured models play the same role inPipeline, the samemodel_namecan be used. When integrating new models, just ensure thatmodel_nameis different from other existing functional models. The corresponding model is fetched throughmodel_namein thePipeline'sfrom_pretrained.model_class: Model architecture import path, pointing to the model architecture class implemented in Step 1, for examplediffsynth.models.qwen_image_text_encoder.QwenImageTextEncoder.state_dict_converter: Optional parameter. If model file format conversion is needed, the import path of the model conversion logic needs to be filled in, for examplediffsynth.utils.state_dict_converters.qwen_image_text_encoder.QwenImageTextEncoderStateDictConverter.extra_kwargs: Optional parameter. If additional parameters need to be passed when initializing the model, these parameters need to be filled in. For example, models DiffSynth-Studio/Qwen-Image-Blockwise-ControlNet-Canny and DiffSynth-Studio/Qwen-Image-Blockwise-ControlNet-Inpaint both adopt theQwenImageBlockWiseControlNetstructure indiffsynth/models/qwen_image_controlnet.py, but the latter also needs additional configurationadditional_in_dim=4. Therefore, this configuration information needs to be filled in theextra_kwargsfield.
We provide a piece of code to quickly understand how models are loaded through this configuration information:
from diffsynth.core import hash_model_file, load_state_dict, skip_model_initialization
from diffsynth.models.qwen_image_text_encoder import QwenImageTextEncoder
from diffsynth.utils.state_dict_converters.qwen_image_text_encoder import QwenImageTextEncoderStateDictConverter
import torch
model_hash = "8004730443f55db63092006dd9f7110e"
model_name = "qwen_image_text_encoder"
model_class = QwenImageTextEncoder
state_dict_converter = QwenImageTextEncoderStateDictConverter
extra_kwargs = {}
model_path = [
"models/Qwen/Qwen-Image/text_encoder/model-00001-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00002-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00003-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00004-of-00004.safetensors",
]
if hash_model_file(model_path) == model_hash:
with skip_model_initialization():
model = model_class(**extra_kwargs)
state_dict = load_state_dict(model_path, torch_dtype=torch.bfloat16, device="cuda")
state_dict = state_dict_converter(state_dict)
model.load_state_dict(state_dict, assign=True)
print("Done!")Q: The logic of the above code looks very simple, why is this part of code in
DiffSynth-Studioextremely complex?A: Because we provide aggressive VRAM management functions that are coupled with the model loading logic, this leads to the complexity of the framework structure. We have tried our best to simplify the interface exposed to developers.
The model_hash in diffsynth/configs/model_configs.py is not uniquely existing. Multiple models may exist in the same model file. For this situation, please use multiple model Configs to load each model separately, and write the corresponding state_dict_converter to separate the parameters required by each model.
Step 4: Verifying Whether the Model Can Be Recognized and Loaded
After model integration, the following code can be used to verify whether the model can be correctly recognized and loaded. The following code will attempt to load the model into memory:
from diffsynth.models.model_loader import ModelPool
model_pool = ModelPool()
model_pool.auto_load_model(
[
"models/Qwen/Qwen-Image/text_encoder/model-00001-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00002-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00003-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00004-of-00004.safetensors",
],
)If the model can be recognized and loaded, you will see the following output:
Loading models from: [
"models/Qwen/Qwen-Image/text_encoder/model-00001-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00002-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00003-of-00004.safetensors",
"models/Qwen/Qwen-Image/text_encoder/model-00004-of-00004.safetensors"
]
Loaded model: {
"model_name": "qwen_image_text_encoder",
"model_class": "diffsynth.models.qwen_image_text_encoder.QwenImageTextEncoder",
"extra_kwargs": null
}Step 5: Writing Model VRAM Management Scheme
DiffSynth-Studio supports complex VRAM management. See Enabling VRAM Management for details.
En Developer Guide Training Diffusion Models (`docs/en/Developer_Guide/Training_Diffusion_Models.md`)
Integrating Model Training
After integrating models and implementing Pipeline, the next step is to integrate model training functionality.
Training-Inference Consistent Pipeline Modification
To ensure strict consistency between training and inference processes, we will use most of the inference code during training, but still need to make minor modifications.
First, add extra logic during inference to switch the image-to-image/video-to-video logic based on the scheduler state. Taking Qwen-Image as an example:
class QwenImageUnit_InputImageEmbedder(PipelineUnit):
def __init__(self):
super().__init__(
input_params=("input_image", "noise", "tiled", "tile_size", "tile_stride"),
output_params=("latents", "input_latents"),
onload_model_names=("vae",)
)
def process(self, pipe: QwenImagePipeline, input_image, noise, tiled, tile_size, tile_stride):
if input_image is None:
return {"latents": noise, "input_latents": None}
pipe.load_models_to_device(['vae'])
image = pipe.preprocess_image(input_image).to(device=pipe.device, dtype=pipe.torch_dtype)
input_latents = pipe.vae.encode(image, tiled=tiled, tile_size=tile_size, tile_stride=tile_stride)
if pipe.scheduler.training:
return {"latents": noise, "input_latents": input_latents}
else:
latents = pipe.scheduler.add_noise(input_latents, noise, timestep=pipe.scheduler.timesteps[0])
return {"latents": latents, "input_latents": input_latents}Then, enable Gradient Checkpointing in model_fn, which will significantly reduce the VRAM required for training at the cost of computational speed. This is not mandatory, but we strongly recommend doing so.
Taking Qwen-Image as an example, before modification:
text, image = block(
image=image,
text=text,
temb=conditioning,
image_rotary_emb=image_rotary_emb,
attention_mask=attention_mask,
)After modification:
from ..core import gradient_checkpoint_forward
text, image = gradient_checkpoint_forward(
block,
use_gradient_checkpointing,
use_gradient_checkpointing_offload,
image=image,
text=text,
temb=conditioning,
image_rotary_emb=image_rotary_emb,
attention_mask=attention_mask,
)Writing Training Scripts
DiffSynth-Studio does not strictly encapsulate the training framework, but exposes the script content to developers. This approach makes it more convenient to modify training scripts to implement additional functions. Developers can refer to existing training scripts, such as examples/qwen_image/model_training/train.py, for modification to adapt to new model training.
En Diffusion Templates Introducing Diffusion Templates (`docs/en/Diffusion_Templates/Introducing_Diffusion_Templates.md`)
Diffusion Templates
Diffusion Templates is a controllable generation plugin framework for Diffusion models in DiffSynth-Studio, providing additional controllable generation capabilities for base models.
- Open-source code: DiffSynth-Studio
- Technical report: arXiv
- Project page: GitHub
- Documentation reference
- Introduction to Diffusion Templates: English Version、中文版
- Detailed Architecture of Diffusion Templates: English Version、中文版
- Template Model Inference: English Version、中文版
- Template Model Training: English Version、中文版
- Online demo: ModelScope
- Model collection: ModelScope、ModelScope International、HuggingFace
| Model Name | ModelScope | ModelScope International | HuggingFace | Inference Code | Low VRAM Inference Code | Training Code | Training Validation Code |
|---|---|---|---|---|---|---|---|
| Structure Control | link | link | link | code | code | code | code |
| Brightness Adjustment | link | link | link | code | code | code | code |
| Color Adjustment | link | link | link | code | code | code | code |
| Image Editing | link | link | link | code | code | code | code |
| Super-Resolution | link | link | link | code | code | code | code |
| Sharpness Enhancement | link | link | link | code | code | code | code |
| Aesthetic Alignment | link | link | link | code | code | code | code |
| Local Redrawing | link | link | link | code | code | code | code |
| Content Reference | link | link | link | code | code | code | code |
| Age Control | link | link | link | code | code | code | code |
| Panda Meme (Easter Egg Model) | link | link | link | code | code | code | code |
| Dataset Name | ModelScope | ModelScope International | HuggingFace |
|---|---|---|---|
| Text-to-Image | link | link | link |
| Local Redrawing | link | link | link |
| Background Replacement | link | link | link |
| Clothing Replacement | link | link | link |
| Pose Adjustment | link | link | link |
| Foreground Modification | link | link | link |
| Local Addition/Removal | link | link | link |
| Super-Resolution | link | link | link |
| Portrait Generation | link | link | link |
| Random Cropping | link | link | link |
| Lighting Adjustment | link | link | link |
| Scene Structure | link | link | link |
| Facial Expression Editing | link | link | link |
| View Angle Adjustment | link | link | link |
| Style Transfer | link | link | link |
| Multi-Resolution | link | link | link |
| Multi-Image Merge | link | link | link |
Model Performance Overview
- Super-Resolution + Sharpness Enhancement: Generate ultra-high-resolution images
| Low Resolution Input | High Resolution Output |
|---|---|
- Structure Control + Aesthetic Alignment + Sharpness Enhancement: Fully-equipped ControlNet
| Structure Control Image | Output Image |
|---|---|
- Structure Control + Image Editing + Color Adjustment: Artistic Style Creation at Will
| Structure Control Image | Editing Input Image | Output Image |
|---|---|---|
- Brightness Control + Image Editing + Local Redrawing: Cross-dimensional Elements in Images
| Reference Image | Redrawing Area | Output Image |
|---|---|---|
En Diffusion Templates Template Model Inference (`docs/en/Diffusion_Templates/Template_Model_Inference.md`)
Template Model Inference
Enabling Template Models on Base Model Pipelines
Using the base model black-forest-labs/FLUX.2-klein-base-4B as an example, when generating images using only the base model:
from diffsynth.diffusion.template import TemplatePipeline
from diffsynth.pipelines.flux2_image import Flux2ImagePipeline, ModelConfig
import torch
# Load base model
pipe = Flux2ImagePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="black-forest-labs/FLUX.2-klein-4B", origin_file_pattern="text_encoder/*.safetensors"),
ModelConfig(model_id="black-forest-labs/FLUX.2-klein-base-4B", origin_file_pattern="transformer/*.safetensors"),
ModelConfig(model_id="black-forest-labs/FLUX.2-klein-4B", origin_file_pattern="vae/diffusion_pytorch_model.safetensors"),
],
tokenizer_config=ModelConfig(model_id="black-forest-labs/FLUX.2-klein-4B", origin_file_pattern="tokenizer/"),
)
# Generate an image
image = pipe(
prompt="a cat",
seed=0, cfg_scale=4,
height=1024, width=1024,
)
image.save("image.png")The Template model DiffSynth-Studio/Template-KleinBase4B-Brightness can control image brightness during generation. Through the TemplatePipeline model, it can be loaded from ModelScope (via ModelConfig(model_id="xxx/xxx")) or from a local path (via ModelConfig(path="xxx")). Inputting scale=0.8 increases image brightness. Note that in the code, input parameters for pipe must be transferred to template_pipeline, and template_inputs should be added.
# Load Template model
template_pipeline = TemplatePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="DiffSynth-Studio/Template-KleinBase4B-Brightness")
],
)
# Generate an image
image = template_pipeline(
pipe,
prompt="a cat",
seed=0, cfg_scale=4,
height=1024, width=1024,
template_inputs=[{"scale": 0.8}],
)
image.save("image_0.8.png")CFG Enhancement for Template Models
Template models can enable CFG (Classifier-Free Guidance) to make control effects more pronounced. For example, with the model DiffSynth-Studio/Template-KleinBase4B-Brightness, adding negative_template_inputs to the TemplatePipeline input parameters and setting its scale to 0.5 will generate images with more noticeable brightness variations by contrasting both sides.
# Generate an image with CFG
image = template_pipeline(
pipe,
prompt="a cat",
seed=0, cfg_scale=4,
height=1024, width=1024,
template_inputs=[{"scale": 0.8}],
negative_template_inputs=[{"scale": 0.5}],
)
image.save("image_0.8_cfg.png")Low VRAM Support
Template models currently do not support the main framework's VRAM management, but lazy loading can be used - loading Template models only when needed for inference. This significantly reduces VRAM requirements when enabling multiple Template models, with peak VRAM usage being that of a single Template model. Add parameter lazy_loading=True to enable.
template_pipeline = TemplatePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="DiffSynth-Studio/Template-KleinBase4B-Brightness")
],
lazy_loading=True,
)The base model's Pipeline and Template Pipeline are completely independent and can enable VRAM management on demand.
When Template model outputs contain LoRA in Template Cache, you need to enable VRAM management for the base model's Pipeline or enable LoRA hot loading (using the code below), otherwise LoRA weights will be fused repeatedly.
pipe.dit = pipe.enable_lora_hot_loading(pipe.dit)Enabling Multiple Template Models
TemplatePipeline can load multiple Template models. During inference, use model_id in template_inputs to distinguish inputs for each Template model.
After enabling VRAM management for the base model's Pipeline and lazy loading for Template Pipeline, you can load any number of Template models.
/* Detailed source-code truncated for AI context efficiency. */Super-Resolution + Sharpness Enhancement
Combining DiffSynth-Studio/Template-KleinBase4B-Upscaler and DiffSynth-Studio/Template-KleinBase4B-Sharpness can upscale blurry images while improving detail clarity.
image = template(
pipe,
prompt="A cat is sitting on a stone.",
seed=0, cfg_scale=4, num_inference_steps=50,
template_inputs = [
{
"model_id": 3,
"image": Image.open("data/examples/templates/image_lowres_100.jpg"),
"prompt": "A cat is sitting on a stone.",
},
{
"model_id": 5,
"scale": 1,
},
],
negative_template_inputs = [
{
"model_id": 3,
"image": Image.open("data/examples/templates/image_lowres_100.jpg"),
"prompt": "",
},
{
"model_id": 5,
"scale": 0,
},
],
)
image.save("image_Upscaler_Sharpness.png")| Low Resolution Input | High Resolution Output |
|---|---|
Structure Control + Aesthetic Alignment + Sharpness Enhancement
DiffSynth-Studio/Template-KleinBase4B-ControlNet controls composition, DiffSynth-Studio/Template-KleinBase4B-Aesthetic fills in details, and DiffSynth-Studio/Template-KleinBase4B-Sharpness ensures clarity. Combining these three Template models produces exquisite images.
image = template(
pipe,
prompt="A cat is sitting on a stone, bathed in bright sunshine.",
seed=0, cfg_scale=4, num_inference_steps=50,
template_inputs = [
{
"model_id": 1,
"image": Image.open("data/examples/templates/image_depth.jpg"),
"prompt": "A cat is sitting on a stone, bathed in bright sunshine.",
},
{
"model_id": 7,
"lora_ids": list(range(1, 180, 2)),
"lora_scales": 2.0,
"merge_type": "mean",
},
{
"model_id": 5,
"scale": 0.8,
},
],
negative_template_inputs = [
{
"model_id": 1,
"image": Image.open("data/examples/templates/image_depth.jpg"),
"prompt": "",
},
{
"model_id": 7,
"lora_ids": list(range(1, 180, 2)),
"lora_scales": 2.0,
"merge_type": "mean",
},
{
"model_id": 5,
"scale": 0,
},
],
)
image.save("image_Controlnet_Aesthetic_Sharpness.png")| Structure Control Image | Output Image |
|---|---|
Structure Control + Image Editing + Color Adjustment
DiffSynth-Studio/Template-KleinBase4B-ControlNet controls composition, DiffSynth-Studio/Template-KleinBase4B-Edit preserves original image details like fur texture, and DiffSynth-Studio/Template-KleinBase4B-SoftRGB controls color tones, creating an artistic masterpiece.
image = template(
pipe,
prompt="A cat is sitting on a stone. Colored ink painting.",
seed=0, cfg_scale=4, num_inference_steps=50,
template_inputs = [
{
"model_id": 1,
"image": Image.open("data/examples/templates/image_depth.jpg"),
"prompt": "A cat is sitting on a stone. Colored ink painting.",
},
{
"model_id": 2,
"image": Image.open("data/examples/templates/image_reference.jpg"),
"prompt": "Convert the image style to colored ink painting.",
},
{
"model_id": 4,
"R": 0.9,
"G": 0.5,
"B": 0.3,
},
],
negative_template_inputs = [
{
"model_id": 1,
"image": Image.open("data/examples/templates/image_depth.jpg"),
"prompt": "",
},
{
"model_id": 2,
"image": Image.open("data/examples/templates/image_reference.jpg"),
"prompt": "",
},
],
)
image.save("image_Controlnet_Edit_SoftRGB.png")| Structure Control Image | Editing Input Image | Output Image |
|---|---|---|
Brightness Control + Image Editing + Local Redrawing
DiffSynth-Studio/Template-KleinBase4B-Brightness generates bright scenes, DiffSynth-Studio/Template-KleinBase4B-Edit references original image layout, and DiffSynth-Studio/Template-KleinBase4B-Inpaint keeps background unchanged, generating cross-dimensional content.
image = template(
pipe,
prompt="A cat is sitting on a stone. Flat anime style.",
seed=0, cfg_scale=4, num_inference_steps=50,
template_inputs = [
{
"model_id": 0,
"scale": 0.6,
},
{
"model_id": 2,
"image": Image.open("data/examples/templates/image_reference.jpg"),
"prompt": "Convert the image style to flat anime style.",
},
{
"model_id": 6,
"image": Image.open("data/examples/templates/image_reference.jpg"),
"mask": Image.open("data/examples/templates/image_mask_1.jpg"),
"force_inpaint": True,
},
],
negative_template_inputs = [
{
"model_id": 0,
"scale": 0.5,
},
{
"model_id": 2,
"image": Image.open("data/examples/templates/image_reference.jpg"),
"prompt": "",
},
{
"model_id": 6,
"image": Image.open("data/examples/templates/image_reference.jpg"),
"mask": Image.open("data/examples/templates/image_mask_1.jpg"),
},
],
)
image.save("image_Brightness_Edit_Inpaint.png")| Reference Image | Redrawing Area | Output Image |
|---|---|---|
En Diffusion Templates Template Model Training (`docs/en/Diffusion_Templates/Template_Model_Training.md`)
Template Model Training
DiffSynth-Studio currently provides comprehensive Template training support for black-forest-labs/FLUX.2-klein-base-4B, with more model adaptations coming soon.
Continuing Training from Pretrained Models
To continue training from our pretrained models, refer to the table in FLUX.2 to find the corresponding training script.
Building New Template Models
Template Model Component Format
A Template model binds to a model repository (or local folder) containing a code file model.py as the entry point. Here's the template for model.py:
import torch
class CustomizedTemplateModel(torch.nn.Module):
def __init__(self):
super().__init__()
@torch.no_grad()
def process_inputs(self, xxx, **kwargs):
yyy = xxx
return {"yyy": yyy}
def forward(self, yyy, **kwargs):
zzz = yyy
return {"zzz": zzz}
class DataProcessor:
def __call__(self, www, **kwargs):
xxx = www
return {"xxx": xxx}
TEMPLATE_MODEL = CustomizedTemplateModel
TEMPLATE_MODEL_PATH = "model.safetensors"
TEMPLATE_DATA_PROCESSOR = DataProcessorDuring Template model inference, Template Input passes through TEMPLATE_MODEL's process_inputs and forward to generate Template Cache.
flowchart LR;
i@{shape: text, label: "Template Input"}-->p[process_inputs];
subgraph TEMPLATE_MODEL
p[process_inputs]-->f[forward]
end
f[forward]-->c@{shape: text, label: "Template Cache"};During Template model training, Template Input comes from the dataset through TEMPLATE_DATA_PROCESSOR.
flowchart LR;
d@{shape: text, label: "Dataset"}-->dp[TEMPLATE_DATA_PROCESSOR]-->p[process_inputs];
subgraph TEMPLATE_MODEL
p[process_inputs]-->f[forward]
end
f[forward]-->c@{shape: text, label: "Template Cache"};`TEMPLATE_MODEL`
TEMPLATE_MODEL implements the Template model logic, inheriting from torch.nn.Module with required process_inputs and forward methods. These two methods form the complete Template model inference process, split into two stages to better support two-stage split training.
process_inputsmust use@torch.no_grad()for gradient-free computationforwardmust contain all gradient computations required for training
Both methods should accept **kwargs for compatibility. Reserved parameters include:
- To interact with the base model Pipeline (e.g., call text encoder), add
pipeparameter to method inputs - To enable Gradient Checkpointing, add
use_gradient_checkpointinganduse_gradient_checkpointing_offloadtoforwardinputs - Multiple Template models use
model_idto distinguish Template Inputs - do not use this field in method parameters
`TEMPLATE_MODEL_PATH` (Optional)
TEMPLATE_MODEL_PATH specifies the relative path to pretrained weights. For example:
TEMPLATE_MODEL_PATH = "model.safetensors"For multi-file models:
TEMPLATE_MODEL_PATH = [
"model-00001-of-00003.safetensors",
"model-00002-of-00003.safetensors",
"model-00003-of-00003.safetensors",
]Set to None for random initialization:
TEMPLATE_MODEL_PATH = None`TEMPLATE_DATA_PROCESSOR` (Optional)
To train Template models with DiffSynth-Studio, datasets should contain template_inputs fields in metadata.json. These fields pass through TEMPLATE_DATA_PROCESSOR to generate inputs for Template model methods.
For example, the brightness control model DiffSynth-Studio/Template-KleinBase4B-Brightness takes scale as input:
[
{
"image": "images/image_1.jpg",
"prompt": "a cat",
"template_inputs": {"scale": 0.2}
},
{
"image": "images/image_2.jpg",
"prompt": "a dog",
"template_inputs": {"scale": 0.6}
}
]class DataProcessor:
def __call__(self, scale, **kwargs):
return {"scale": scale}
TEMPLATE_DATA_PROCESSOR = DataProcessorOr calculate scale from image paths:
[
{
"image": "images/image_1.jpg",
"prompt": "a cat",
"template_inputs": {"image": "/path/to/your/dataset/images/image_1.jpg"}
}
]class DataProcessor:
def __call__(self, image, **kwargs):
image = Image.open(image)
image = np.array(image)
return {"scale": image.astype(np.float32).mean() / 255}
TEMPLATE_DATA_PROCESSOR = DataProcessorTraining Template Models
A Template model is "trainable" if its Template Cache variables are fully decoupled from the base model Pipeline - these variables should reach model_fn without participating in any Pipeline Unit calculations.
For training with black-forest-labs/FLUX.2-klein-base-4B, use these training script parameters:
--extra_inputs: Additional inputs. Usetemplate_inputsfor text-to-image models,edit_image,template_inputsfor image editing models--template_model_id_or_path: Template model ID or local path (use:suffix for ModelScope IDs, e.g.,"DiffSynth-Studio/Template-KleinBase4B-Brightness:")--remove_prefix_in_ckpt: State dict prefix to remove when saving models (use"pipe.template_model.")--trainable_models: Trainable components (use"template_model"for full model, or"template_model.xxx,template_model.yyy"for specific components)
Example training script:
accelerate launch examples/flux2/model_training/train.py \
--dataset_base_path data/diffsynth_example_dataset/flux2/Template-KleinBase4B-Brightness \
--dataset_metadata_path data/diffsynth_example_dataset/flux2/Template-KleinBase4B-Brightness/metadata.jsonl \
--extra_inputs "template_inputs" \
--max_pixels 1048576 \
--dataset_repeat 50 \
--model_id_with_origin_paths "black-forest-labs/FLUX.2-klein-4B:text_encoder/*.safetensors,black-forest-labs/FLUX.2-klein-base-4B:transformer/*.safetensors,black-forest-labs/FLUX.2-klein-4B:vae/diffusion_pytorch_model.safetensors" \
--template_model_id_or_path "examples/flux2/model_training/scripts/brightness" \
--tokenizer_path "black-forest-labs/FLUX.2-klein-4B:tokenizer/" \
--learning_rate 1e-4 \
--num_epochs 2 \
--remove_prefix_in_ckpt "pipe.template_model." \
--output_path "./models/train/Template-KleinBase4B-Brightness_example" \
--trainable_models "template_model" \
--use_gradient_checkpointing \
--find_unused_parametersInteracting with Base Model Pipeline Components
Template models can interact with base model Pipelines. For example, using the text encoder:
class CustomizedTemplateModel(torch.nn.Module):
def __init__(self):
super().__init__()
self.xxx = xxx()
@torch.no_grad()
def process_inputs(self, text, pipe, **kwargs):
input_ids = pipe.tokenizer(text)
text_emb = pipe.text_encoder(input_ids)
return {"text_emb": text_emb}
def forward(self, text_emb, pipe, **kwargs):
kv_cache = self.xxx(text_emb)
return {"kv_cache": kv_cache}
TEMPLATE_MODEL = CustomizedTemplateModelUsing Non-Trainable Components
For models with pretrained components:
class CustomizedTemplateModel(torch.nn.Module):
def __init__(self):
super().__init__()
self.image_encoder = XXXEncoder.from_pretrained(xxx)
self.mlp = MLP()
@torch.no_grad()
def process_inputs(self, image, **kwargs):
emb = self.image_encoder(image)
return {"emb": emb}
def forward(self, emb, **kwargs):
kv_cache = self.mlp(emb)
return {"kv_cache": kv_cache}
TEMPLATE_MODEL = CustomizedTemplateModelSet --trainable_models template_model.mlp to train only the MLP component.
Training on Low VRAM Devices
The framework supports splitting Template model training into two stages: the first stage performs gradient-free computation, and the second stage performs gradient updates. For more information, refer to the documentation: Two-stage Split Training. Here's a sample script:
/* Detailed source-code truncated for AI context efficiency. */Two-stage split training can reduce VRAM requirements and improve training speed. The training process is lossless in precision, but requires significant disk space for storing cache files.
To further reduce VRAM requirements, you can enable fp8 precision by adding the parameters --fp8_models "black-forest-labs/FLUX.2-klein-4B:text_encoder/*.safetensors,black-forest-labs/FLUX.2-klein-4B:vae/diffusion_pytorch_model.safetensors" and --fp8_models "black-forest-labs/FLUX.2-klein-base-4B:transformer/*.safetensors" to the two-stage training. Note that fp8 precision can only be enabled on non-trainable model components and introduces minor errors.
Uploading Template Models
After training, follow these steps to upload Template models to ModelScope for wider distribution.
- Set model path in
model.py:
TEMPLATE_MODEL_PATH = "model.safetensors"- Upload using ModelScope CLI:
modelscope upload user_name/your_model_id /path/to/your/model.py model.py --token ms-xxx- Package model files:
from diffsynth.diffusion.template import load_template_model, load_state_dict
from safetensors.torch import save_file
import torch
model = load_template_model("path/to/your/template/model", torch_dtype=torch.bfloat16, device="cpu")
state_dict = load_state_dict("path/to/your/ckpt/epoch-1.safetensors", torch_dtype=torch.bfloat16, device="cpu")
state_dict.update(model.state_dict())
save_file(state_dict, "model.safetensors")- Upload model file:
modelscope upload user_name/your_model_id /path/to/your/model/epoch-1.safetensors model.safetensors --token ms-xxx- Verify inference:
from diffsynth.diffusion.template import TemplatePipeline
from diffsynth.pipelines.flux2_image import Flux2ImagePipeline, ModelConfig
import torch
# Load base model
pipe = Flux2ImagePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="black-forest-labs/FLUX.2-klein-4B", origin_file_pattern="text_encoder/*.safetensors"),
ModelConfig(model_id="black-forest-labs/FLUX.2-klein-base-4B", origin_file_pattern="transformer/*.safetensors"),
ModelConfig(model_id="black-forest-labs/FLUX.2-klein-4B", origin_file_pattern="vae/diffusion_pytorch_model.safetensors"),
],
tokenizer_config=ModelConfig(model_id="black-forest-labs/FLUX.2-klein-4B", origin_file_pattern="tokenizer/"),
)
# Load Template model
template_pipeline = TemplatePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="user_name/your_model_id")
],
)
# Generate image
image = template_pipeline(
pipe,
prompt="a cat",
seed=0, cfg_scale=4,
height=1024, width=1024,
template_inputs=[{xxx}],
)
image.save("image.png")
---
### En Diffusion Templates Understanding Diffusion Templates (`docs/en/Diffusion_Templates/Understanding_Diffusion_Templates.md`)
# Diffusion Templates Architecture Details
The Diffusion Templates framework is a controllable generation plugin framework in DiffSynth-Studio that provides additional controllable generation capabilities for Diffusion models.
## Framework Structure
The Diffusion Templates framework structure is shown below:
```mermaid
flowchart TD;
subgraph Template Pipeline
si@{shape: text, label: "Template Input"}-->i1@{shape: text, label: "Template Input 1"};
si@{shape: text, label: "Template Input"}-->i2@{shape: text, label: "Template Input 2"};
si@{shape: text, label: "Template Input"}-->i3@{shape: text, label: "Template Input 3"};
i1@{shape: text, label: "Template Input 1"}-->m1[Template Model 1]-->c1@{shape: text, label: "Template Cache 1"};
i2@{shape: text, label: "Template Input 2"}-->m2[Template Model 2]-->c2@{shape: text, label: "Template Cache 2"};
i3@{shape: text, label: "Template Input 3"}-->m3[Template Model 3]-->c3@{shape: text, label: "Template Cache 3"};
c1-->c@{shape: text, label: "Template Cache"};
c2-->c;
c3-->c;
end
i@{shape: text, label: "Model Input"}-->m[Diffusion Pipeline]-->o@{shape: text, label: "Model Output"};
c-->m;The framework contains these module designs:
- Template Input: Template model input. Format: Python dictionary with fields determined by each Template model (e.g.,
{"scale": 0.8}) - Template Model: Template model, loadable from ModelScope (
ModelConfig(model_id="xxx/xxx")) or local path (ModelConfig(path="xxx")) - Template Cache: Template model output. Format: Python dictionary with fields matching base model Pipeline input parameters
- Template Pipeline: Module for managing multiple Template models. Handles model loading and cache integration
When the Diffusion Templates framework is disabled, base model components (Text Encoder, DiT, VAE) are loaded into the Diffusion Pipeline. Model Input (prompt, height, width) produces Model Output (e.g., images).
When enabled, Template models are loaded into the Template Pipeline. The Template Pipeline outputs Template Cache (a subset of Diffusion Pipeline input parameters) for subsequent processing in the Diffusion Pipeline. This enables controllable generation by intercepting part of the Diffusion Pipeline's input parameters.
Model Capability Medium
Template Cache is defined as a subset of Diffusion Pipeline input parameters, ensuring framework generality. We restrict Template model inputs to only be Diffusion Pipeline parameters. The KV-Cache is particularly suitable as a Diffusion medium:
- Proven effective in LLM Skills (prompts are converted to KV-Cache)
- Has "high permission" in Diffusion models - can directly control image generation
- Supports sequence-level concatenation for multiple Template models
- Requires minimal development (add pipeline parameter and integrate to model)
Other potential Template mediums:
- Residual: Used in ControlNet for point-to-point control, but has resolution limitations and potential conflicts when merging
- LoRA: Treated as input parameters rather than model components
Currently, we only support KV-Cache and LoRA as Template Cache mediums in FLUX.2 Pipeline, with plans to support more models and mediums in the future.
Template Model Format
A Template model has this structure:
Template_Model
├── model.py
└── model.safetensorsWhere model.py is the entry point and model.safetensors contains model weights. For implementation details, see Template Model Training or existing Template models.
En Model Details ACE Step (`docs/en/Model_Details/ACE-Step.md`)
ACE-Step
ACE-Step 1.5 is an open-source music generation model based on DiT architecture, supporting text-to-music, audio cover, repainting and other functionalities, running efficiently on consumer-grade hardware.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the ACE-Step/Ace-Step1.5 model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 3GB VRAM.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| ACE-Step/Ace-Step1.5 | code | code | code | code | code | code |
| ACE-Step/acestep-v15-turbo-shift1 | code | code | code | code | code | code |
| ACE-Step/acestep-v15-turbo-shift3 | code | code | code | code | code | code |
| ACE-Step/acestep-v15-turbo-continuous | code | code | code | code | code | code |
| ACE-Step/acestep-v15-base | code | code | code | code | code | code |
| ACE-Step/acestep-v15-base: CoverTask | code | code | — | — | — | — |
| ACE-Step/acestep-v15-base: RepaintTask | code | code | — | — | — | — |
| ACE-Step/acestep-v15-sft | code | code | code | code | code | code |
| ACE-Step/acestep-v15-xl-base | code | code | code | code | code | code |
| ACE-Step/acestep-v15-xl-sft | code | code | code | code | code | code |
| ACE-Step/acestep-v15-xl-turbo | code | code | code | code | code | code |
| DiffSynth-Studio/acestep15xlsft-lora-music | code | code | code | code | - | - |
Model Inference
The model is loaded via AceStepPipeline.from_pretrained, see Loading Models for details.
The input parameters for AceStepPipeline inference include:
prompt: Text description of the music.cfg_scale: Classifier-free guidance scale, defaults to 1.0.lyrics: Lyrics text.task_type: Task type,可选 values include"text2music"(text-to-music),"cover"(audio cover),"repaint"(repainting), defaults to"text2music".reference_audios: List of reference audio tensors for timbre reference.src_audio: Source audio tensor for cover or repaint tasks.denoising_strength: Denoising strength, controlling how much the output is influenced by source audio, defaults to 1.0.audio_cover_strength: Audio cover step ratio, controlling how many steps use cover condition in cover tasks, defaults to 1.0.audio_code_string: Input audio code string for cover tasks with discrete audio codes.repainting_ranges: List of repainting time ranges (tuples of floats, in seconds) for repaint tasks.repainting_strength: Repainting intensity, controlling the degree of change in repainted areas, defaults to 1.0.duration: Audio duration in seconds, defaults to 60.bpm: Beats per minute, defaults to 100.keyscale: Musical key scale, defaults to "B minor".timesignature: Time signature, defaults to "4".vocal_language: Vocal language, defaults to "unknown".seed: Random seed.rand_device: Device for noise generation, defaults to "cpu".num_inference_steps: Number of inference steps, defaults to 8.shift: Timestep shift parameter for the scheduler, defaults to 3.0.
Model Training
Models in the ace_step series are trained uniformly via examples/ace_step/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- ACE-Step Specific Parameters
--tokenizer_path: Tokenizer path, in format model_id:origin_pattern.--silence_latent_path: Silence latent path, in format model_id:origin_pattern.--initialize_model_on_cpu: Whether to initialize models on CPU.
Example Dataset
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details Anima (`docs/en/Model_Details/Anima.md`)
Anima
Anima is an image generation model trained and open-sourced by CircleStone Labs and Comfy Org.
Installation
Before using this project for model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more installation information, please refer to Install Dependencies.
Quick Start
The following code demonstrates how to quickly load the circlestone-labs/Anima model for inference. VRAM management is enabled by default, allowing the framework to automatically control model parameter loading based on available VRAM. Minimum 8GB VRAM required.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Validation after Full Training | LoRA Training | Validation after LoRA Training |
|---|---|---|---|---|---|---|
| circlestone-labs/Anima | code | code | code | code | code | code |
Special training scripts:
- Differential LoRA Training: doc
- FP8 Precision Training: doc
- Two-Stage Split Training: doc
- End-to-End Direct Distillation: doc
Model Inference
Models are loaded through AnimaImagePipeline.from_pretrained, see Model Inference for details.
Input parameters for AnimaImagePipeline inference include:
prompt: Text description of the desired image content.negative_prompt: Content to exclude from the generated image (default:"").cfg_scale: Classifier-free guidance parameter (default: 4.0).input_image: Input image for image-to-image generation (default:None).denoising_strength: Controls similarity to input image (default: 1.0).height: Image height (must be multiple of 16, default: 1024).width: Image width (must be multiple of 16, default: 1024).seed: Random seed (default:None).rand_device: Device for random noise generation (default:"cpu").num_inference_steps: Inference steps (default: 30).sigma_shift: Scheduler sigma offset (default:None).progress_bar_cmd: Progress bar implementation (default:tqdm.tqdm).
For VRAM constraints, enable VRAM Management. Recommended low-VRAM configurations are provided in the "Model Overview" table above.
Model Training
Anima models are trained through examples/anima/model_training/train.py with parameters including:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Dataset root directory.--dataset_metadata_path: Metadata file path.--dataset_repeat: Dataset repetition per epoch.--dataset_num_workers: Dataloader worker count.--data_file_keys: Metadata fields to load (comma-separated).
- Model Loading
--model_paths: Model paths (JSON format).--model_id_with_origin_paths: Model IDs with origin paths (e.g.,"anima-team/anima-1B:text_encoder/*.safetensors").--extra_inputs: Additional pipeline inputs (e.g.,controlnet_inputsfor ControlNet).--fp8_models: FP8-formatted models (same format as--model_paths).--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Training Configuration
--learning_rate: Learning rate.--num_epochs: Training epochs.--trainable_models: Trainable components (e.g.,dit,vae,text_encoder).--find_unused_parameters: Handle unused parameters in DDP training.--weight_decay: Weight decay value.--task: Training task (default:sft).
- Output Configuration
--output_path: Model output directory.--remove_prefix_in_ckpt: Remove state dict prefixes.--save_steps: Model saving interval.
- LoRA Configuration
--lora_base_model: Target model for LoRA.--lora_target_modules: Target modules for LoRA.--lora_rank: LoRA rank.--lora_checkpoint: LoRA checkpoint path.--preset_lora_path: Preloaded LoRA checkpoint path.--preset_lora_model: Model to merge LoRA with (e.g.,dit).
- Gradient Configuration
--use_gradient_checkpointing: Enable gradient checkpointing.--use_gradient_checkpointing_offload: Offload checkpointing to CPU.--gradient_accumulation_steps: Gradient accumulation steps.
- Image Resolution
--height: Image height (empty for dynamic resolution).--width: Image width (empty for dynamic resolution).--max_pixels: Maximum pixel area for dynamic resolution.
- Dataset Configuration
- Anima-Specific Parameters
--tokenizer_path: Tokenizer path for text-to-image models.--tokenizer_t5xxl_path: T5-XXL tokenizer path.
We provide a sample image dataset for testing:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetFor training script details, refer to Model Training. For advanced training techniques, see Training Framework Documentation.
En Model Details Boogu Image (`docs/en/Model_Details/Boogu-Image.md`)
Boogu-Image
Boogu-Image supports text-to-image, image-to-image, and instruction-guided image editing.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the Boogu/Boogu-Image-0.1-Base model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 8GB VRAM.
from diffsynth.pipelines.boogu_image import BooguImagePipeline, ModelConfig
import torch
vram_config = {
"offload_dtype": torch.float8_e4m3fn,
"offload_device": "cpu",
"onload_dtype": torch.float8_e4m3fn,
"onload_device": "cpu",
"preparing_dtype": torch.float8_e4m3fn,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
}
pipe = BooguImagePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="Boogu/Boogu-Image-0.1-Base", origin_file_pattern="transformer/*.safetensors", **vram_config),
ModelConfig(model_id="Boogu/Boogu-Image-0.1-Base", origin_file_pattern="mllm/*.safetensors", **vram_config),
ModelConfig(model_id="Boogu/Boogu-Image-0.1-Base", origin_file_pattern="vae/*.safetensors", **vram_config),
],
processor_config=ModelConfig(model_id="Boogu/Boogu-Image-0.1-Base", origin_file_pattern="mllm/"),
vram_limit=torch.cuda.mem_get_info("cuda")[1] / (1024 ** 3) - 0.5,
)
output = pipe(
prompt="a cat",
negative_prompt="",
height=1024,
width=1024,
seed=42,
num_inference_steps=50,
cfg_scale=4.0,
)
output.save("image_Boogu-Image-0.1-Base.jpg")Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| Boogu/Boogu-Image-0.1-Base | code | code | code | code | code | code |
| Boogu/Boogu-Image-0.1-Turbo | code | code | code | code | code | code |
| Boogu/Boogu-Image-0.1-Edit | code | code | code | code | code | code |
Model Inference
The model is loaded via BooguImagePipeline.from_pretrained, see Loading Models for details.
The input parameters for BooguImagePipeline inference include:
prompt: Text prompt describing the desired content or editing instruction.negative_prompt: Negative prompt specifying what should not appear in the result, defaults to empty string.cfg_scale: Classifier-free guidance scale factor, defaults to 4.0. Higher values make the output more closely follow the prompt.input_image: Input image for image-to-image (img2img). When provided, the input image is noised and denoised according todenoising_strength.edit_image: Image to be edited for instruction-guided editing. When provided, the model modifies the image according to thepromptinstruction.height: Height of the output image, defaults to 1024. Must be divisible by 16.width: Width of the output image, defaults to 1024. Must be divisible by 16.seed: Random seed for reproducibility. Set toNonefor random seed.denoising_strength: Denoising strength controlling how much the input image is repainted, defaults to 1.0. Only effective wheninput_imageis provided.sigmas: Custom sigma scheduling sequence to override the default scheduling strategy. Required for Turbo models.num_inference_steps: Number of inference steps, defaults to 20. More steps typically yield better quality.max_sequence_length: Maximum sequence length for the text encoder, defaults to 1280.max_input_image_pixels: Maximum pixel area for input images, defaults to 4194304. Images larger than this will be scaled down.max_input_image_side_length: Maximum side length for input images, defaults to 4096.max_vlm_input_pil_pixels: Maximum pixel area for VLM input images, defaults to 147456. Only effective in image editing mode.max_vlm_input_pil_side_length: Maximum side length for VLM input images, defaults to 768. Only effective in image editing mode.rand_device: Device for generating initial noise, defaults to "cpu".progress_bar_cmd: Progress bar display mode, defaults to tqdm.
When running low on VRAM, please refer to VRAM Management to enable VRAM management features.
Model Training
Models in the boogu_image series are trained uniformly via examples/boogu_image/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- Boogu-Image Specific Parameters
--processor_path: Path to the processor for processing text and image encoder inputs.--initialize_model_on_cpu: Whether to initialize models on CPU. By default, models are initialized on the accelerator device.
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details DiffSynth Music (`docs/en/Model_Details/DiffSynth-Music.md`)
DiffSynth-Music
DiffSynth-Music is a controllable music generation model suite developed by the DiffSynth-Studio team, built upon ACE-Step-1.5. It retains the foundational architecture of ACE-Step-1.5 while incorporating several additional modules to support diverse control mechanisms, including beat, vocals, accompaniment, rhythm, and timbre references.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the DiffSynth-Studio/DiffSynth-Music model for inference. The example covers six use cases: native music synthesis, beats control, vocals control, accompaniment control, prosody control, and audio reference (timbre) control.
/* Detailed source-code truncated for AI context efficiency. */A low VRAM version is available at examples/diffsynth_music/model_inference_low_vram/DiffSynth-Music.py.
Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| DiffSynth-Studio/DiffSynth-Music | code | code | code | code | - | - |
Model Inference
The model is loaded via DiffSynthMusicPipeline.from_pretrained, see Loading Models for details.
The input parameters for DiffSynthMusicPipeline inference include:
prompt: The music description, specifying genre, mood and style.negative_prompt: The negative prompt for classifier-free guidance.lyrics: The lyrics. Structure tags such as[Intro]and[Verse]must each be on their own line. Leave it empty to generate instrumental music.bpm: Beats per minute. Default: 100.timesignature: Time signature. Default: "4".keyscale: Key and scale. Default: "B minor".input_audio: Optional audio tensor used as the starting point, combined withdenoising_strengthfor audio-to-audio generation.duration: Duration of the generated audio in seconds.num_inference_steps: Number of flow-matching steps.cfg_scale: Classifier-free guidance scale.seed: Random seed.tiled: Whether to use tiled VAE decoding.kv_cache: Pre-computed KV cache for faster inference.target_audio/target_track: Optional track fusion, replacing the given track of the output with the corresponding track oftarget_audio.
The TemplatePipeline provides three control adapters loaded via template_inputs with model_id:
model_id: 0— template_control: beats / vocals / accompaniment control. The generated audio follows the rhythm of the given audio. Combined withtarget_audio, the target track is fused into the output.model_id: 1— template_prosody: the generated vocals follow the pitch and rhythm (prosody) of the given audio. Useextract_prosodyto build the input.model_id: 2— template_reference: the timbre of the output is controlled by the reference audio.
If you run out of VRAM, please refer to VRAM Management.
Model Training
The model supports LoRA training and full training of the DiT (trainable models: dit) on (prompt, lyrics, audio) data. See examples/diffsynth_music/model_training for training scripts.
En Model Details ERNIE Image (`docs/en/Model_Details/ERNIE-Image.md`)
ERNIE-Image
ERNIE-Image is a powerful image generation model with 8B parameters developed by Baidu, featuring a compact and efficient architecture with strong instruction-following capability. Based on an 8B DiT backbone, it delivers performance comparable to larger (20B+) models in certain scenarios while maintaining parameter efficiency. It offers reliable performance in instruction understanding and execution, text generation (English/Chinese/Japanese), and overall stability.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the PaddlePaddle/ERNIE-Image model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 3G VRAM.
from diffsynth.pipelines.ernie_image import ErnieImagePipeline, ModelConfig
import torch
vram_config = {
"offload_dtype": torch.bfloat16,
"offload_device": "cpu",
"onload_dtype": torch.bfloat16,
"onload_device": "cpu",
"preparing_dtype": torch.bfloat16,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
}
pipe = ErnieImagePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device='cuda',
model_configs=[
ModelConfig(model_id="PaddlePaddle/ERNIE-Image", origin_file_pattern="transformer/diffusion_pytorch_model*.safetensors", **vram_config),
ModelConfig(model_id="PaddlePaddle/ERNIE-Image", origin_file_pattern="text_encoder/model.safetensors", **vram_config),
ModelConfig(model_id="PaddlePaddle/ERNIE-Image", origin_file_pattern="vae/diffusion_pytorch_model.safetensors", **vram_config),
],
tokenizer_config=ModelConfig(model_id="PaddlePaddle/ERNIE-Image", origin_file_pattern="tokenizer/"),
vram_limit=torch.cuda.mem_get_info("cuda")[1] / (1024 ** 3) - 0.5,
)
image = pipe(
prompt="一只黑白相间的中华田园犬",
negative_prompt="",
height=1024,
width=1024,
seed=42,
num_inference_steps=50,
cfg_scale=4.0,
)
image.save("output.jpg")Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| PaddlePaddle/ERNIE-Image | code | code | code | code | code | code |
| PaddlePaddle/ERNIE-Image-Turbo | code | code | — | — | — | — |
Model Inference
The model is loaded via ErnieImagePipeline.from_pretrained, see Loading Models for details.
The input parameters for ErnieImagePipeline inference include:
prompt: The prompt describing the content to appear in the image.negative_prompt: The negative prompt describing what should not appear in the image, default value is"".cfg_scale: Classifier-free guidance parameter, default value is 4.0.height: Image height, must be a multiple of 16, default value is 1024.width: Image width, must be a multiple of 16, default value is 1024.seed: Random seed. Default isNone, meaning completely random.rand_device: The computing device for generating random Gaussian noise matrices, default is"cuda". When set tocuda, different GPUs will produce different results.num_inference_steps: Number of inference steps, default value is 50.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low-VRAM configurations for each model in the "Model Overview" table above.
Model Training
ERNIE-Image series models are trained uniformly via examples/ernie_image/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, e.g.,"PaddlePaddle/ERNIE-Image:transformer/diffusion_pytorch_model*.safetensors", separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image. Leave empty to enable dynamic resolution.--width: Width of the image. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.
- Dataset Configuration
- ERNIE-Image Specific Parameters
--tokenizer_path: Path to the tokenizer, leave empty to auto-download from remote.
We provide an example image dataset for testing, which can be downloaded with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details FLUX.3 (`docs/en/Model_Details/FLUX.3.md`)
FLUX.3
FLUX.3 Image Klein is a 5.6B open-weight image generation model from Black Forest Labs that supports text-to-image generation, editing with up to 10 reference images, and bounding-box layout control in one model.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Model Access
black-forest-labs/flux-3-image-klein and black-forest-labs/flux-3-image-klein-base are access-request-only models on ModelScope. Before downloading them:
- Request access on the model page (Klein or Klein Base).
- Log in via the CLI with your ModelScope token:
modelscope login --token ms-xxx. Get your token at Settings / Access Token.
Access must be granted for the same account used to log in.
Quick Start
Running the following code will load the black-forest-labs/flux-3-image-klein model for inference. VRAM management is enabled, and the framework automatically controls parameter loading based on available VRAM. Minimum 4GB VRAM is required to run.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| black-forest-labs/flux-3-image-klein | code | code | code | code | code | code |
| black-forest-labs/flux-3-image-klein-base | code | code | code | code | code | code |
| DiffSynth-Studio/flux-3-image-klein-entropack | code | code | - | - | code | code |
Model Inference
The model is loaded via Flux3ImagePipeline.from_pretrained, see Loading Models for details.
The input parameters for Flux3ImagePipeline inference include:
prompt: Prompt describing the target image, edit, or layout; defaults to"".negative_prompt: Negative prompt used by Base when CFG is enabled; defaults to"".cfg_scale: Classifier-free guidance strength, defaults to 1.0. Distilled requires 1.0; the Base examples use 3.5.input_image: Input image for image-to-image generation; generation starts from pure noise when omitted.denoising_strength: Image-to-image denoising strength from 0 to 1, defaults to 1.0.edit_image: Reference image(s) for editing, accepting onePIL.Imageor a list of images; up to 10 references are supported.height: Output image height, defaults to 1024 and is aligned to a multiple of 16.width: Output image width, defaults to 1024 and is aligned to a multiple of 16.seed: Random seed, defaults toNone.rand_device: Device used to generate the initial Gaussian noise, defaults to"cpu".num_inference_steps: Number of inference steps, defaults to 16. The Distilled examples use 16 steps and the Base examples use 50.alpha: Rational-shift coefficient for the FLUX.3 flow-matching schedule, defaults to 6.93.progress_bar_cmd: Function wrapping the inference timesteps with a progress bar, defaults totqdm.
If VRAM is insufficient, please enable VRAM Management. Recommended configurations are linked in the "Model Overview" table above.
Model Training
Models in the FLUX.3 series are trained uniformly via examples/flux3/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- FLUX.3 Specific Parameters
--processor_path: Local processor path or model ID; when omitted, the Base model'stext_encoder/is used.--initialize_model_on_cpu: Initialize models on CPU to reduce peak VRAM during distributed training startup.
We have built a sample image dataset for your testing. You can download this dataset with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details FLUX (`docs/en/Model_Details/FLUX.md`)
FLUX
FLUX is an image generation model series developed and open-sourced by Black Forest Labs.
Installation
Before using this project for model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information about installation, please refer to Install Dependencies.
Quick Start
Run the following code to quickly load the black-forest-labs/FLUX.1-dev model and perform inference. VRAM management is enabled, and the framework will automatically control model parameter loading based on remaining VRAM. Minimum 8GB VRAM is required to run.
import torch
from diffsynth.pipelines.flux_image import FluxImagePipeline, ModelConfig
vram_config = {
"offload_dtype": torch.float8_e4m3fn,
"offload_device": "cpu",
"onload_dtype": torch.float8_e4m3fn,
"onload_device": "cpu",
"preparing_dtype": torch.float8_e4m3fn,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
}
pipe = FluxImagePipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="black-forest-labs/FLUX.1-dev", origin_file_pattern="flux1-dev.safetensors", **vram_config),
ModelConfig(model_id="black-forest-labs/FLUX.1-dev", origin_file_pattern="text_encoder/model.safetensors", **vram_config),
ModelConfig(model_id="black-forest-labs/FLUX.1-dev", origin_file_pattern="text_encoder_2/*.safetensors", **vram_config),
ModelConfig(model_id="black-forest-labs/FLUX.1-dev", origin_file_pattern="ae.safetensors", **vram_config),
],
vram_limit=torch.cuda.mem_get_info("cuda")[1] / (1024 ** 3) - 1,
)
prompt = "CG, masterpiece, best quality, solo, long hair, wavy hair, silver hair, blue eyes, blue dress, medium breasts, dress, underwater, air bubble, floating hair, refraction, portrait. The girl's flowing silver hair shimmers with every color of the rainbow and cascades down, merging with the floating flora around her."
image = pipe(prompt=prompt, seed=0)
image.save("image.jpg")Model Overview
Special Training Scripts:
- Differential LoRA Training: doc
- FP8 Precision Training: doc
- Two-stage Split Training: doc
- End-to-end Direct Distillation: doc
Model Inference
Models are loaded via FluxImagePipeline.from_pretrained, see Loading Models.
Input parameters for FluxImagePipeline inference include:
prompt: Prompt describing the content appearing in the image.negative_prompt: Negative prompt describing content that should not appear in the image, default value is"".cfg_scale: Classifier-free guidance parameter, default value is 1. When set to a value greater than 1, CFG is enabled.height: Image height, must be a multiple of 16.width: Image width, must be a multiple of 16.seed: Random seed. Default isNone, meaning completely random.rand_device: Computing device for generating random Gaussian noise matrix, default is"cpu". When set tocuda, different GPUs will produce different generation results.num_inference_steps: Number of inference steps, default value is 30.embedded_guidance: Embedded guidance parameter, default value is 3.5.t5_sequence_length: Sequence length of the T5 text encoder, default is 512.tiled: Whether to enable VAE tiling inference, default isFalse. Setting toTruecan significantly reduce VRAM usage during VAE encoding/decoding stages, producing slight errors and slightly longer inference time.tile_size: Tile size during VAE encoding/decoding stages, default is 128, only effective whentiled=True.tile_stride: Tile stride during VAE encoding/decoding stages, default is 64, only effective whentiled=True, must be less than or equal totile_size.progress_bar_cmd: Progress bar, default istqdm.tqdm. Can be disabled by setting tolambda x:x.controlnet_inputs: ControlNet model inputs, type isControlNetInputlist.ipadapter_images: IP-Adapter model input image list.ipadapter_scale: Guidance strength of the IP-Adapter model.infinityou_id_image: InfiniteYou model input image.infinityou_guidance: Guidance strength of the InfiniteYou model.kontext_images: Kontext model input images.eligen_entity_prompts: EliGen partition control prompt list.eligen_entity_masks: EliGen partition control region mask image list.eligen_enable_on_negative: Whether to enable EliGen partition control on the negative side of CFG.eligen_enable_inpaint: Whether to enable EliGen partition control inpainting function.lora_encoder_inputs: LoRA encoder input image list.lora_encoder_scale: Guidance strength of the LoRA encoder.step1x_reference_image: Step1X model reference image.flex_inpaint_image: Flex model image to be inpainted.flex_inpaint_mask: Flex model inpainting mask.flex_control_image: Flex model control image.flex_control_strength: Flex model control strength.flex_control_stop: Flex model control stop timestep.nexus_gen_reference_image: Nexus-Gen model reference image.flux_fill_image: FLUX.1-Fill model image to be inpainted.flux_fill_mask: FLUX.1-Fill model inpainting mask.flux_redux_image: FLUX.1-Redux model reference image.insert_anything_source_image: Insert-Anything model source image, i.e., the target image to be edited.insert_anything_source_mask: Insert-Anything model source image mask, specifying the region to be edited.insert_anything_ref_image: Insert-Anything model reference image, providing the content to be inserted.insert_anything_ref_mask: Insert-Anything model reference image mask, specifying the target object in the reference image.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low VRAM configurations for each model in the example code, see the table in the "Model Overview" section above.
Model Training
FLUX series models are uniformly trained through examples/flux/model_training/train.py, and the script parameters include:
- General Training Parameters
- Dataset Basic Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Metadata file path of the dataset.--dataset_repeat: Number of times the dataset is repeated in each epoch.--dataset_num_workers: Number of processes for each DataLoader.--data_file_keys: Field names to be loaded from metadata, usually image or video file paths, separated by,.
- Model Loading Configuration
--model_paths: Paths of models to be loaded. JSON format.--model_id_with_origin_paths: Model IDs with original paths, e.g.,"black-forest-labs/FLUX.1-dev:flux1-dev.safetensors". Separated by commas.--extra_inputs: Extra input parameters required by the model Pipeline, e.g.,controlnet_inputswhen training ControlNet models, separated by,.--fp8_models: Models loaded in FP8 format, consistent with--model_pathsor--model_id_with_origin_pathsformat. Currently only supports models whose parameters are not updated by gradients (no gradient backpropagation, or gradients only update their LoRA).--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Training Basic Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether there are unused parameters in DDP training. Some models contain redundant parameters that do not participate in gradient calculation, and this setting needs to be enabled to avoid errors in multi-GPU training.--weight_decay: Weight decay size, see torch.optim.AdamW.--task: Training task, default issft. Some models support more training modes, please refer to the documentation of each specific model.
- Output Configuration
--output_path: Model saving path.--remove_prefix_in_ckpt: Remove prefix in the state dict of the model file.--save_steps: Interval of training steps to save the model. If this parameter is left blank, the model is saved once per epoch.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path of the LoRA checkpoint. If this path is provided, LoRA will be loaded from this checkpoint.--preset_lora_path: Preset LoRA checkpoint path. If this path is provided, this LoRA will be loaded in the form of being merged into the base model. This parameter is used for LoRA differential training.--preset_lora_model: Model that the preset LoRA is merged into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Image Width/Height Configuration (Applicable to Image Generation and Video Generation Models)
--height: Height of image or video. Leaveheightandwidthblank to enable dynamic resolution.--width: Width of image or video. Leaveheightandwidthblank to enable dynamic resolution.--max_pixels: Maximum pixel area of image or video frames. When dynamic resolution is enabled, images with resolution larger than this value will be downscaled, and images with resolution smaller than this value will remain unchanged.
- Dataset Basic Configuration
- FLUX Specific Parameters
--tokenizer_1_path: Path of the CLIP tokenizer, leave blank to automatically download from remote.--tokenizer_2_path: Path of the T5 tokenizer, leave blank to automatically download from remote.--align_to_opensource_format: Whether to align LoRA format to open-source format, only applicable to DiT's LoRA.
We have built a sample image dataset for your testing. You can download this dataset with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe have written recommended training scripts for each model, please refer to the table in the "Model Overview" section above. For how to write model training scripts, please refer to Model Training; for more advanced training algorithms, please refer to Training Framework Detailed Explanation.
En Model Details FLUX2 (`docs/en/Model_Details/FLUX2.md`)
FLUX.2
FLUX.2 is an image generation model trained and open-sourced by Black Forest Labs.
Installation
Before using this project for model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information about installation, please refer to Install Dependencies.
Quick Start
Run the following code to quickly load the black-forest-labs/FLUX.2-dev model and perform inference. VRAM management is enabled, and the framework will automatically control model parameter loading based on remaining VRAM. Minimum 10GB VRAM is required to run.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
Special Training Scripts:
- Differential LoRA Training: doc
- FP8 Precision Training: doc
- Two-stage Split Training: doc
- End-to-end Direct Distillation: doc
Model Inference
Models are loaded via Flux2ImagePipeline.from_pretrained, see Loading Models.
Input parameters for Flux2ImagePipeline inference include:
prompt: Prompt describing the content appearing in the image.negative_prompt: Negative prompt describing content that should not appear in the image, default value is"".cfg_scale: Classifier-free guidance parameter, default value is 1. When set to a value greater than 1, CFG is enabled.height: Image height, must be a multiple of 16.width: Image width, must be a multiple of 16.seed: Random seed. Default isNone, meaning completely random.rand_device: Computing device for generating random Gaussian noise matrix, default is"cpu". When set tocuda, different GPUs will produce different generation results.num_inference_steps: Number of inference steps, default value is 30.embedded_guidance: Embedded guidance parameter, default value is 3.5.t5_sequence_length: Sequence length of the T5 text encoder, default is 512.tiled: Whether to enable VAE tiling inference, default isFalse. Setting toTruecan significantly reduce VRAM usage during VAE encoding/decoding stages, producing slight errors and slightly longer inference time.tile_size: Tile size during VAE encoding/decoding stages, default is 128, only effective whentiled=True.tile_stride: Tile stride during VAE encoding/decoding stages, default is 64, only effective whentiled=True, must be less than or equal totile_size.progress_bar_cmd: Progress bar, default istqdm.tqdm. Can be disabled by setting tolambda x:x.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low VRAM configurations for each model in the example code, see the table in the "Model Overview" section above.
Model Training
FLUX.2 series models are uniformly trained through examples/flux2/model_training/train.py, and the script parameters include:
- General Training Parameters
- Dataset Basic Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Metadata file path of the dataset.--dataset_repeat: Number of times the dataset is repeated in each epoch.--dataset_num_workers: Number of processes for each DataLoader.--data_file_keys: Field names to be loaded from metadata, usually image or video file paths, separated by,.
- Model Loading Configuration
--model_paths: Paths of models to be loaded. JSON format.--model_id_with_origin_paths: Model IDs with original paths, e.g.,"black-forest-labs/FLUX.2-dev:text_encoder/*.safetensors". Separated by commas.--extra_inputs: Extra input parameters required by the model Pipeline, e.g.,controlnet_inputswhen training ControlNet models, separated by,.--fp8_models: Models loaded in FP8 format, consistent with--model_pathsor--model_id_with_origin_pathsformat. Currently only supports models whose parameters are not updated by gradients (no gradient backpropagation, or gradients only update their LoRA).--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Training Basic Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether there are unused parameters in DDP training. Some models contain redundant parameters that do not participate in gradient calculation, and this setting needs to be enabled to avoid errors in multi-GPU training.--weight_decay: Weight decay size, see torch.optim.AdamW.--task: Training task, default issft. Some models support more training modes, please refer to the documentation of each specific model.
- Output Configuration
--output_path: Model saving path.--remove_prefix_in_ckpt: Remove prefix in the state dict of the model file.--save_steps: Interval of training steps to save the model. If this parameter is left blank, the model is saved once per epoch.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path of the LoRA checkpoint. If this path is provided, LoRA will be loaded from this checkpoint.--preset_lora_path: Preset LoRA checkpoint path. If this path is provided, this LoRA will be loaded in the form of being merged into the base model. This parameter is used for LoRA differential training.--preset_lora_model: Model that the preset LoRA is merged into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Image Width/Height Configuration (Applicable to Image Generation and Video Generation Models)
--height: Height of image or video. Leaveheightandwidthblank to enable dynamic resolution.--width: Width of image or video. Leaveheightandwidthblank to enable dynamic resolution.--max_pixels: Maximum pixel area of image or video frames. When dynamic resolution is enabled, images with resolution larger than this value will be downscaled, and images with resolution smaller than this value will remain unchanged.
- Dataset Basic Configuration
- FLUX.2 Specific Parameters
--tokenizer_path: Path of the tokenizer, applicable to text-to-image models, leave blank to automatically download from remote.
We have built a sample image dataset for your testing. You can download this dataset with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe have written recommended training scripts for each model, please refer to the table in the "Model Overview" section above. For how to write model training scripts, please refer to Model Training; for more advanced training algorithms, please refer to Training Framework Detailed Explanation.
En Model Details Fibo (`docs/en/Model_Details/Fibo.md`)
Fibo 2.0
Fibo 2.0 is an 8B unified image generation and editing model developed by Bria AI. It runs in the latent space of the FLUX.2 VAE, encodes prompts and reference images with Qwen3-VL-4B, and uses a single-stream DiT with 44 blocks. Text-to-image and image editing share one set of weights.
Fibo 2.0 Dev is the full, non-distilled base model, trained with classifier-free guidance and intended for fine-tuning. Fibo 2.0 Edit Turbo is a distilled editing-only variant for 4-step, guidance-free inference.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Run the following code to quickly load the briaai/Fibo-2.0-dev model and perform text-to-image inference. VRAM management is enabled, and the framework will automatically control model parameter loading based on remaining VRAM.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| briaai/Fibo-2.0-dev | code | code | code | code | code | code |
| briaai/Fibo-2.0-edit-turbo | code | code | - | - | - | - |
The training columns above link to the text-to-image scripts. Fibo 2.0 Dev also supports image-editing training, see Model Training.
Model Inference
The model is loaded via FiboPipeline.from_pretrained, see Loading Models for details.
from_pretrained also accepts:
processor_config:ModelConfigof the Qwen3-VL processor.font_config:ModelConfigof the fonts used to drawtext_renderimages. It can point to a font directory, a font file, or a list of font files.
The input parameters for FiboPipeline inference include:
prompt: A structured JSON caption, a plain description, or an edit instruction. Edit prompts may also be a JSON caption of the result with anedit_instructionfield.negative_prompt: Used whencfg_scale > 1. Defaults to the null caption Fibo 2.0 was trained with.edit_image: The image or list of reference images to edit. Leave empty for text-to-image.mask: For single-image editing, white marks the region to regenerate.text_images: Images of the text in a caption'stext_renderentries. When omitted, they are drawn from the caption.fonts: Font paths used to draw characters that Pillow's built-in font lacks.render_text_images: Whether to draw text images from the caption, default isTrue.height,width: Output size, multiples of 16. When editing and left empty, the trained size for the first image's aspect ratio is used.cfg_scale: Classifier-free guidance scale, default is 5.0.num_inference_steps: Number of inference steps, default is 30.denoising_strength: Denoising strength, range is 0~1, default is 1. Lower values start from a lower noise level.seed,rand_device: Random seed and the device used to generate noise.progress_bar_cmd: Progress bar, default istqdm.tqdm. Set tolambda x:xto disable it.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low VRAM configurations for each model in the example code, see the table in the "Model Overview" section above.
Image Editing
edit_image accepts one image, or a list of up to 5 images for multi-reference editing. The example below uses the distilled Fibo-2.0-edit-turbo model:
from PIL import Image
from diffsynth.pipelines.fibo import FiboPipeline, ModelConfig
pipe = FiboPipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="briaai/Fibo-2.0-edit-turbo", origin_file_pattern="transformer/*.safetensors"),
ModelConfig(model_id="briaai/Fibo-2.0-edit-turbo", origin_file_pattern="text_encoder/model.safetensors"),
ModelConfig(model_id="briaai/Fibo-2.0-edit-turbo", origin_file_pattern="vae/diffusion_pytorch_model.safetensors"),
],
processor_config=ModelConfig(model_id="briaai/Fibo-2.0-edit-turbo", origin_file_pattern="processor/"),
font_config=ModelConfig(model_id="briaai/Fibo-2.0-edit-turbo", origin_file_pattern="fonts/"),
)
source = Image.open("photo.jpg").convert("RGB")
image = pipe(prompt="replace the jacket with a red leather one", edit_image=source, seed=42, num_inference_steps=4, cfg_scale=1.0)
image.save("edited.jpg")Fibo-2.0-dev uses 50 steps with cfg_scale=5.0; Fibo-2.0-edit-turbo uses 4 steps with cfg_scale=1.0.
Masked Editing
For single-image editing, pass mask together with edit_image (white marks the region to regenerate):
mask = Image.open("mask.png").convert("L") # white: regenerate, black: keep
image = pipe(
prompt="replace the suitcase with a vintage brown leather one",
edit_image=source, mask=mask, seed=42, num_inference_steps=4, cfg_scale=1.0,
)Text Rendering
Fibo renders text from a caption's text_render entries. Each entry's text is drawn as a small black-on-white image and shown to Qwen3-VL together with the prompt:
caption = {
"short_description": "A young woman holding a white sign that reads \"ModelScope\".",
"objects": [{"description": "a white sign with bold black text", "location": "center"}],
"style_medium": "photograph",
"text_render": [{"text": "ModelScope"}],
}
image = pipe(
prompt=json.dumps(caption, separators=(",", ":"), ensure_ascii=False),
height=1024, width=1024, seed=42, num_inference_steps=50, cfg_scale=5.0,
)font_config loads the fonts shipped with the model. To render characters that those fonts lack, pass additional font paths through fonts:
image = pipe(prompt=prompt, fonts=["/path/to/a/cjk-font.ttf"], ...)Characters not covered by any font are drawn with Pillow's built-in font. Text rendering accuracy is model-dependent, especially for non-Latin scripts.
Model Training
Fibo 2.0 Dev is trained via examples/fibo/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Basic Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Metadata file path of the dataset.--dataset_repeat: Number of times the dataset is repeated in each epoch.--dataset_num_workers: Number of processes for each DataLoader.--data_file_keys: Field names to be loaded from metadata, usually image or video file paths, separated by,.
- Model Loading Configuration
--model_paths: Paths of models to be loaded. JSON format.--model_id_with_origin_paths: Model IDs with original paths, e.g.,"briaai/Fibo-2.0-dev:transformer/*.safetensors". Separated by commas.--extra_inputs: Extra input parameters required by the model Pipeline, e.g.,edit_imageandmaskwhen training the image editing model, separated by,.--fp8_models: Models loaded in FP8 format, consistent with--model_pathsor--model_id_with_origin_pathsformat. Currently only supports models whose parameters are not updated by gradients (no gradient backpropagation, or gradients only update their LoRA).--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Training Basic Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether there are unused parameters in DDP training. Some models contain redundant parameters that do not participate in gradient calculation, and this setting needs to be enabled to avoid errors in multi-GPU training.--weight_decay: Weight decay size, see torch.optim.AdamW.--task: Training task, default issft. Some models support more training modes, please refer to the documentation of each specific model.
- Output Configuration
--output_path: Model saving path.--remove_prefix_in_ckpt: Remove prefix in the state dict of the model file.--save_steps: Interval of training steps to save the model. If this parameter is left blank, the model is saved once per epoch.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path of the LoRA checkpoint. If this path is provided, LoRA will be loaded from this checkpoint.--preset_lora_path: Preset LoRA checkpoint path. If this path is provided, this LoRA will be loaded in the form of being merged into the base model. This parameter is used for LoRA differential training.--preset_lora_model: Model that the preset LoRA is merged into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Image Width/Height Configuration (Applicable to Image Generation and Video Generation Models)
--height: Height of image or video. Leaveheightandwidthblank to enable dynamic resolution.--width: Width of image or video. Leaveheightandwidthblank to enable dynamic resolution.--max_pixels: Maximum pixel area of image or video frames. When dynamic resolution is enabled, images with resolution larger than this value will be downscaled, and images with resolution smaller than this value will remain unchanged.
- Dataset Basic Configuration
- Fibo 2.0 Specific Parameters
--processor_path: Path ormodel_id:origin_file_patternof the Qwen3-VL processor, leave blank to automatically download from remote.--font_path: Path ormodel_id:origin_file_patternof the fonts used for text rendering, leave blank to automatically download from remote.--initialize_model_on_cpu: Whether to initialize the model on CPU.
We have built sample datasets, including a text-to-image dataset and an image-editing dataset. You can download them with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --include "fibo/*" --local_dir ./data/diffsynth_example_datasetdata/diffsynth_example_dataset/fibo/Fibo-2.0-dev: text-to-image, withmetadata.csvcontainingimageandprompt.data/diffsynth_example_dataset/fibo/Fibo-2.0-dev-Edit: image editing, withmetadata.jsoncontainingimage,prompt,edit_image, and an optionalmask.
Fibo 2.0 Dev supports both text-to-image and image-editing training. Training reuses the inference Pipeline: the VAE and text encoder are frozen, and the DiT is trained with flow matching. Full training uses --trainable_models "dit", while LoRA training uses --lora_base_model "dit". For image editing, edit_image and mask are loaded with UnifiedDataset and injected through --extra_inputs.
| Mode | Full Training | LoRA Training | Full Validation | LoRA Validation |
|---|---|---|---|---|
| Text-to-image | script | script | script | script |
| Image editing | script | script | script | script |
The full training scripts use --config_file examples/fibo/model_training/full/accelerate_config.yaml; the LoRA scripts run on a single GPU by default. mask is optional per sample, and rows without mask are trained as ordinary editing samples. Fibo 2.0 Edit Turbo is not intended for fine-tuning; train on Fibo 2.0 Dev instead.
We have written recommended training scripts for each model, please refer to the tables above. For how to write model training scripts, please refer to Model Training; for more advanced training algorithms, please refer to Training Framework Detailed Explanation.
License
Fibo 2.0 is released under CC BY-NC 4.0 for non-commercial use. Commercial use requires a separate agreement with Bria AI.
En Model Details HiDream O1 Image (`docs/en/Model_Details/HiDream-O1-Image.md`)
HiDream-O1-Image
HiDream-O1-Image is an image generation model open-sourced by HiDream.ai, based on the Pixel-Level Unified Transformer (UiT) architecture. This model unifies VAE, DiT, and TextEncoder within a single Qwen3VLModel, performing diffusion denoising directly in pixel patch space without requiring a separate VAE component.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will quickly load the HiDream-ai/HiDream-O1-Image model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 3GB VRAM.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| HiDream-ai/HiDream-O1-Image | code | code | code | code | code | code |
| HiDream-ai/HiDream-O1-Image-Dev | code | code | code | code | code | code |
| DiffSynth-Studio/HidreamO1-i2L-v2 | code | code | code | code | - | - |
Model Inference
The model is loaded via HiDreamO1ImagePipeline.from_pretrained, see Loading Models for details.
The input parameters for HiDreamO1ImagePipeline inference include:
prompt: Text prompt.negative_prompt: Negative prompt, defaults to" ".cfg_scale: Classifier-Free Guidance scale, defaults to 4.0. For the Dev model, it is recommended to set to 1.0.height: Output image height, defaults to 2048.width: Output image width, defaults to 2048.seed: Random seed, defaults to random.rand_device: Noise generation device, defaults to"cpu".num_inference_steps: Number of inference steps, defaults to 50 for Full model and 28 for Dev model.model_type: Model type,"full"for Full model,"dev"for distilled Dev model.shift: Timestep shift parameter affecting sigma computation, defaults to 3.0.noise_scale: Noise scaling factor, defaults to 8.0. For the Dev model, it is recommended to set to 7.5.edit_image: List of reference images for image editing. Defaults to None (text-to-image mode).keep_original_aspect: Whether to preserve the original aspect ratio of reference images, defaults to True.
VRAM Note: HiDream-O1-Image has a large parameter count (~8B). When generating 2048x2048 images, it is recommended to enable VRAM management (vram_config) or use the low VRAM inference scripts.
Model Training
Models in the hidream_o1_image series are trained uniformly via examples/hidream_o1_image/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- HiDream-O1-Image Specific Parameters
--processor_config: Path to the processor configuration file, used for loading AutoProcessor for text tokenization.--noise_scale: Noise scaling factor, defaults to 8.0.--initialize_model_on_cpu: Whether to initialize the model on CPU, which can help reduce peak GPU VRAM usage.
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details Ideogram 4 (`docs/en/Model_Details/Ideogram-4.md`)
Ideogram 4
Ideogram 4 is an image generation model open-sourced by Ideogram. DiffSynth-Studio supports inference, low VRAM inference, full training, and LoRA training for both the FP8 quantized version and the BF16 repackaged version.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information about installation, please refer to Install Dependencies.
Quick Start
Running the following code will load the ideogram-ai/ideogram-4-fp8 model for inference. A minimum of 24GB VRAM is required to run.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| ideogram-ai/ideogram-4-fp8 | code | - | - | - | - | - |
| DiffSynth-Studio/ideogram-4-bf16-repackage | code | code | code | - | code | code |
Model Inference
The model is loaded via Ideogram4Pipeline.from_pretrained, see Loading Models for details.
The input parameters for Ideogram4Pipeline inference include:
prompt: Prompt describing the content appearing in the image. Ideogram 4 supports structured JSON format prompts, including high-level description, style description, and compositional deconstruction.negative_prompt: Negative prompt describing content that should not appear in the image, default value is"".cfg_scale: Classifier-free guidance parameter, default value is 7.0.input_image: Input image for image-to-image generation, used in conjunction withdenoising_strength.denoising_strength: Denoising strength, range is 0~1, default value is 1. When the value approaches 0, the generated image is similar to the input image; when the value approaches 1, the generated image differs more from the input image. Wheninput_imageparameter is not provided, do not set this to a non-1 value.height: Image height, must be a multiple of 16, default value is 1024.width: Image width, must be a multiple of 16, default value is 1024.seed: Random seed. Default isNone, meaning completely random.rand_device: Computing device for generating random Gaussian noise matrix, default is"cpu".num_inference_steps: Number of inference steps, default value is 50.
Model Training
Models in the ideogram4 series are trained uniformly via examples/ideogram4/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- Ideogram-4 Specific Parameters
--tokenizer_path: Path to tokenizer. Defaults to downloading fromideogram-ai/ideogram-4-fp8.
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details Image Quality Metrics (`docs/en/Model_Details/Image-Quality-Metrics.md`)
Image Quality Evaluation Metrics
DiffSynth-Studio provides a suite of image quality evaluation metrics and reward models in diffsynth.metrics to assess text alignment, aesthetic quality, human preference, and image distribution quality of generated images. Example code for these metrics can be found in examples/image_quality_metric/.
Installation
Before using this project for model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information about installation, please refer to Install Dependencies.
Quick Start
Run the following code to quickly load PickScore and score an image against a prompt. The default models will be downloaded from ModelScope to ./models.
from diffsynth.metrics import PickScoreMetric, ModelConfig
from modelscope import dataset_snapshot_download
from PIL import Image
dataset_snapshot_download(
"DiffSynth-Studio/diffsynth_example_dataset",
allow_file_pattern="flux/FLUX.1-dev/*",
local_dir="./data/diffsynth_example_dataset",
)
image = Image.open("data/diffsynth_example_dataset/flux/FLUX.1-dev/1.jpg").convert("RGB")
prompt = "a dog"
metric = PickScoreMetric.from_pretrained(
model_config=ModelConfig(model_id="DiffSynth-Studio/ImageMetrics", origin_file_pattern="PickScore/model.safetensors"),
device="cuda"
)
score = metric.compute(prompt, image)[0]
print(f"PickScore score:: {score:.3f}")Metrics Overview
| Metric | Input | Output | Example Code |
|---|---|---|---|
| PickScore | prompt + PIL Image | Preference Score | code |
| ImageReward | prompt + PIL Image | Preference Score | code |
| HPSv2 | prompt + PIL Image | Preference Score | code |
| HPSv3 | prompt + PIL Image | Preference Score | code |
| CLIP Score | prompt + PIL Image | Text-Image Similarity | code |
| UnifiedReward 2.0 | prompt + PIL Image | multi-dimension scores | code |
| Qwen-Image-Bench | prompt + PIL Image | Overall score and multi-level dimension scores | code |
| UnifiedReward Edit | editing instruction + source image + edited image | Image editing quality score | code |
| Aesthetic | PIL Image | Aesthetic Score | code |
| FID | reference image directory + generated image directory | Distribution Distance | code |
Text-Image Alignment and Preference Evaluation
Applicable metrics: PickScore, ImageReward, HPSv2, HPSv3, CLIP Score, UnifiedReward 2.0, Qwen-Image-Bench
These models are used to evaluate whether an image follows the prompt and aligns with human visual preferences. They must receive both the prompt and the image simultaneously.
Basic Scoring
score = metric.compute(prompt, image)[0]Batch Scoring
If you need to evaluate multiple images, you can directly pass a list:
scores = metric.compute("a cute cat", [image1, image2, image3])
scores = metric.compute(["a cat", "a dog"], [image_cat, image_dog])When prompt is a single string, the same prompt will be applied to every image. When prompt is a list of strings, the number of prompts must exactly match the number of images.
Multi-Dimensional Image Quality Evaluation
Applicable metrics: UnifiedReward 2.0, Qwen-Image-Bench
These metrics also receive a prompt and an image, but in addition to the primary score, evaluate() returns more detailed evaluation dimensions. They are useful when you need to analyze text-image alignment, visual coherence, style, or multi-level quality dimensions.
Qwen-Image-Bench
from diffsynth.metrics import ModelConfig, QwenImageBenchMetric
metric = QwenImageBenchMetric.from_pretrained(
model_config=ModelConfig(
model_id="Qwen/Qwen-Image-Bench",
origin_file_pattern="model-*.safetensors",
),
processor_config=ModelConfig(
model_id="Qwen/Qwen-Image-Bench",
origin_file_pattern="",
),
device="cuda",
)
details = metric.evaluate(prompt, image)[0]
score = details["total_score"]
print(details["level1_scores"])
print(details["level2_scores"])If you only need the primary score, you can also call metric.compute(prompt, image).
Image Editing Quality Evaluation
Applicable metric: UnifiedReward Edit
UnifiedReward Edit evaluates whether an edited image follows the editing instruction and whether it is over-edited. The input usually includes an editing instruction, a source image, and edited image candidates. It supports three tasks:
edit_pointwise_score: scores a single edited result with[source_image, edited_image].edit_pairwise_rank: compares two edited results and returns the winner with[source_image, edited_image_1, edited_image_2].edit_pairwise_score: returns separate scores for two edited results with[source_image, edited_image_1, edited_image_2].
from diffsynth.metrics import ModelConfig, UnifiedRewardEditMetric
metric = UnifiedRewardEditMetric.from_pretrained(
model_config=ModelConfig(
model_id="DiffSynth-Studio/ImageMetrics",
origin_file_pattern="UnifiedReward-Edit-qwen3vl-8b/model-*.safetensors",
),
processor_config=ModelConfig(
model_id="DiffSynth-Studio/ImageMetrics",
origin_file_pattern="UnifiedReward-Edit-qwen3vl-8b/",
),
device="cuda",
)
details = metric.evaluate(
instruction,
[source_image, edited_image],
task="edit_pointwise_score",
)[0]
print(details["score"], details["editing_success"], details["overediting"])Pure Image Aesthetics Evaluation
Applicable metric: Aesthetic
This model solely evaluates aesthetic features such as the composition, color, and clarity of the image itself. It does not require a prompt.
from diffsynth.metrics import AestheticMetric
metric = AestheticMetric.from_pretrained(device="cuda")
score = metric.compute(image)[0]Dataset Distribution Evaluation
Applicable metric: FID (Fréchet Inception Distance)
FID does not score individual images; instead, it compares the overall feature distribution distance between a real reference image set and a generated image set. A lower score indicates that the generated distribution is closer to the real distribution.
from diffsynth.metrics import FIDMetric
reference_dir = "path/to/real_reference_images"
generated_dir = "path/to/model_generated_images"
metric = FIDMetric.from_pretrained(device="cuda", batch_size=16)
fid_score = metric.compute(reference_dir, generated_dir)
print(f"FID: {fid_score:.3f}")The baseline for FID is not fixed or unique. For general image generation, COCO Validation is commonly used; for specific domains (such as medical images or e-commerce products), a reference_dir composed of real data from that specific domain should be provided.
Important Notes
- The scores from PickScore, ImageReward, HPSv2, HPSv3, CLIPScore, UnifiedReward 2.0, Qwen-Image-Bench, UnifiedReward Edit, and Aesthetic are suitable for relative comparison within the same metric. It is not recommended to directly compare the numerical values across different metrics.
- HPSv3, UnifiedReward 2.0, UnifiedReward Edit, and Qwen-Image-Bench are based on multimodal large models, requiring significantly more VRAM than CLIP-based metrics.
- FID is sensitive to the choice of reference, the reference sample size, and the generated sample size.
En Model Details JoyAI Image (`docs/en/Model_Details/JoyAI-Image.md`)
JoyAI-Image
JoyAI-Image is a unified multi-modal foundation model open-sourced by JD.com, supporting image understanding, text-to-image generation, and instruction-guided image editing.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the jd-opensource/JoyAI-Image-Edit model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 4GB VRAM.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| jd-opensource/JoyAI-Image-Edit | code | code | code | code | code | code |
Model Inference
The model is loaded via JoyAIImagePipeline.from_pretrained, see Loading Models for details.
The input parameters for JoyAIImagePipeline inference include:
prompt: Text prompt describing the desired image editing effect.negative_prompt: Negative prompt specifying what should not appear in the result, defaults to empty string.cfg_scale: Classifier-free guidance scale factor, defaults to 5.0. Higher values make the output more closely follow the prompt.edit_image: Image to be edited.denoising_strength: Denoising strength controlling how much the input image is repainted, defaults to 1.0.height: Height of the output image, defaults to 1024. Must be divisible by 16.width: Width of the output image, defaults to 1024. Must be divisible by 16.seed: Random seed for reproducibility. Set toNonefor random seed.max_sequence_length: Maximum sequence length for the text encoder, defaults to 4096.num_inference_steps: Number of inference steps, defaults to 30. More steps typically yield better quality.tiled: Whether to enable tiling for reduced VRAM usage, defaults to False.tile_size: Tile size, defaults to (30, 52).tile_stride: Tile stride, defaults to (15, 26).shift: Shift parameter for the scheduler, controlling the Flow Match scheduling curve, defaults to 4.0.progress_bar_cmd: Progress bar display mode, defaults to tqdm.
Model Training
Models in the joyai_image series are trained uniformly via examples/joyai_image/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image/video. Leave empty to enable dynamic resolution.--width: Width of the image/video. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.--num_frames: Number of frames for video (video generation models only).
- Dataset Configuration
- JoyAI-Image Specific Parameters
--processor_path: Path to the processor for processing text and image encoder inputs.--initialize_model_on_cpu: Whether to initialize models on CPU. By default, models are initialized on the accelerator device.
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
En Model Details Krea 2 (`docs/en/Model_Details/Krea-2.md`)
Krea-2
Krea-2 is an image generation model developed by the Krea team.
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the krea/Krea-2-Raw model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 24GB VRAM.
from diffsynth.pipelines.krea2 import Krea2Pipeline, ModelConfig
import torch
vram_config = {
"offload_dtype": "disk",
"offload_device": "disk",
"onload_dtype": torch.float8_e4m3fn,
"onload_device": "cpu",
"preparing_dtype": torch.float8_e4m3fn,
"preparing_device": "cuda",
"computation_dtype": torch.bfloat16,
"computation_device": "cuda",
}
pipe = Krea2Pipeline.from_pretrained(
torch_dtype=torch.bfloat16,
device="cuda",
model_configs=[
ModelConfig(model_id="krea/Krea-2-Raw", origin_file_pattern="raw.safetensors", **vram_config),
ModelConfig(model_id="Qwen/Qwen3-VL-4B-Instruct", origin_file_pattern="*.safetensors", **vram_config),
ModelConfig(model_id="Qwen/Qwen-Image", origin_file_pattern="vae/diffusion_pytorch_model.safetensors", **vram_config),
],
tokenizer_config=ModelConfig(model_id="Qwen/Qwen3-VL-4B-Instruct", origin_file_pattern=""),
vram_limit=torch.cuda.mem_get_info("cuda")[1] / (1024 ** 3) - 1,
)
prompt = "A cat standing on a stone."
image = pipe(prompt, seed=0, num_inference_steps=52, cfg_scale=4.5)
image.save("image.jpg")Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| krea/Krea-2-Raw | code | code | code | code | code | code |
| krea/Krea-2-Turbo | code | code | code | code | code | code |
Model Inference
The model is loaded via Krea2Pipeline.from_pretrained, see Loading Models for details.
The input parameters for Krea2Pipeline inference include:
prompt: Prompt describing the content of the image to generate, default value is"".negative_prompt: Negative prompt describing content that should not appear in the image, default value is"".cfg_scale: Classifier-free guidance parameter, default value is 3.5.height: Image height, must be a multiple of 16, default value is 1024.width: Image width, must be a multiple of 16, default value is 1024.seed: Random seed, default isNone, meaning completely random.rand_device: Computing device for generating random Gaussian noise matrix, default is"cpu".num_inference_steps: Number of inference steps, default value is 52.mu: Timestep dynamic shift parameter, default isNone.progress_bar_cmd: Progress bar, default istqdm.tqdm. Can be disabled by setting tolambda x:x.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low VRAM configurations for each model in the example code, see the table in the "Model Overview" section above.
Model Training
Models in the Krea-2 series are trained uniformly via examples/krea2/model_training/train.py. The script parameters include:
- General Training Parameters
- Dataset Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Path to the dataset metadata file.--dataset_repeat: Number of dataset repeats per epoch.--dataset_num_workers: Number of processes per DataLoader.--data_file_keys: Field names to load from metadata, typically paths to image or video files, separated by,.
- Model Loading Configuration
--model_paths: Paths to load models from, in JSON format.--model_id_with_origin_paths: Model IDs with original paths, separated by commas.--extra_inputs: Additional input parameters required by the model Pipeline, separated by,.--fp8_models: Models to load in FP8 format, currently only supported for models whose parameters are not updated by gradients.--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Basic Training Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether unused parameters exist in DDP training.--weight_decay: Weight decay magnitude.--task: Training task, defaults tosft.
- Output Configuration
--output_path: Path to save the model.--remove_prefix_in_ckpt: Remove prefix in the model's state dict.--save_steps: Interval in training steps to save the model.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path to LoRA checkpoint.--preset_lora_path: Path to preset LoRA checkpoint for LoRA differential training.--preset_lora_model: Which model to integrate preset LoRA into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to CPU memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Resolution Configuration
--height: Height of the image. Leave empty to enable dynamic resolution.--width: Width of the image. Leave empty to enable dynamic resolution.--max_pixels: Maximum pixel area, images larger than this will be scaled down during dynamic resolution.
- Dataset Configuration
- Krea-2 Specific Parameters
--tokenizer_path: Path to the tokenizer, leave blank to automatically download from remote.--initialize_model_on_cpu: Whether to initialize models on CPU.--align_to_opensource_format: Whether to align the LoRA format to the opensource format, useful for compatibility with other frameworks.
We have built a sample dataset for your testing. You can download it with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --include "krea2/*" --local_dir ./data/diffsynth_example_datasetWe provide recommended training scripts for each model, please refer to the table in "Model Overview" above. For guidance on writing model training scripts, see Model Training; for more advanced training algorithms, see Training Framework Overview.
License
⚠️ Notice: Krea-2 weights (Raw and Turbo) are released under the Krea 2 Community License, not the Apache 2.0 license that governs DiffSynth-Studio itself.
En Model Details LTX 2 (`docs/en/Model_Details/LTX-2.md`)
LTX-2
LTX-2 is a series of audio-video generation models developed by Lightricks.
Installation
Before using this project for model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information about installation, please refer to Installation Dependencies.
Quick Start
Run the following code to quickly load the Lightricks/LTX-2.3 model and perform inference. VRAM management has been enabled, and the framework will automatically control model parameter loading based on remaining VRAM. It can run with a minimum of 8GB VRAM.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
Model Inference
Models are loaded through LTX2AudioVideoPipeline.from_pretrained, see Loading Models for details.
Input parameters for LTX2AudioVideoPipeline inference include:
prompt: Prompt describing the content appearing in the video.negative_prompt: Negative prompt describing content that should not appear in the video, default value is"".denoising_strength: Denoising strength, range is 0~1, default value is 1.0.input_images: List of input images for image-to-video generation.input_images_indexes: Frame index list of input images in the video.input_images_strength: Strength of input images, default value is 1.0.seed: Random seed. Default isNone, which means completely random.rand_device: Computing device for generating random Gaussian noise matrix, default is"cpu". When set tocuda, different results will be generated on different GPUs.height: Video height, must be a multiple of 32 (single-stage) or 64 (two-stage).width: Video width, must be a multiple of 32 (single-stage) or 64 (two-stage).num_frames: Number of video frames, default value is 121, must be a multiple of 8 + 1.auto_duration: Predict the clip duration from the prompt. Defaults toFalse. When enabled,num_framesis not required and the Duration Head must be loaded.auto_duration_min_seconds/auto_duration_max_seconds: Lower and upper bounds (seconds) for the predicted duration. Default to 1.0 and 20.0.audio_only: Whether to generate audio only (T2A). Defaults toFalse; whenTrue, the video VAE and latent upsampler are not required.cfg_scale: Classifier-free guidance parameter, default value is 3.0.num_inference_steps: Number of inference steps, default value is 40.tiled: Whether to enable VAE tiling inference, default isTrue. When set toTrue, it can significantly reduce VRAM usage during VAE encoding/decoding stages, with slight errors and minor inference time extension.tile_size_in_pixels: Pixel tiling size during VAE encoding/decoding stages, default is 512.tile_overlap_in_pixels: Pixel tiling overlap size during VAE encoding/decoding stages, default is 128.tile_size_in_frames: Frame tiling size during VAE encoding/decoding stages, default is 128.tile_overlap_in_frames: Frame tiling overlap size during VAE encoding/decoding stages, default is 24.use_two_stage_pipeline: Whether to use two-stage pipeline, default isFalse.use_distilled_pipeline: Whether to use distilled pipeline, default isFalse.progress_bar_cmd: Progress bar, default istqdm.tqdm. Can be set tolambda x:xto hide the progress bar.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low VRAM configurations for each model in the example code, see the table in the previous "Supported Inference Scripts" section.
Model Training
LTX-2 series models are uniformly trained through examples/ltx2/model_training/train.py, and the script parameters include:
- General Training Parameters
- Dataset Basic Configuration
--dataset_base_path: Root directory of the dataset.--dataset_metadata_path: Metadata file path of the dataset.--dataset_repeat: Number of times the dataset is repeated in each epoch.--dataset_num_workers: Number of processes for each DataLoader.--data_file_keys: Field names to be loaded from metadata, usually image or video file paths, separated by,.
- Model Loading Configuration
--model_paths: Paths of models to be loaded. JSON format.--model_id_with_origin_paths: Model IDs with original paths, e.g.,"Wan-AI/Wan2.1-T2V-1.3B:diffusion_pytorch_model*.safetensors". Separated by commas.--extra_inputs: Extra input parameters required by the model Pipeline, e.g., extra parameters when training image editing models, separated by,.--fp8_models: Models loaded in FP8 format, consistent with--model_pathsor--model_id_with_origin_pathsformat. Currently only supports models whose parameters are not updated by gradients (no gradient backpropagation, or gradients only update their LoRA).--quant_options: Dynamically quantize loaded models. Semicolon-separated entries, each<model_string>:<method>[/<exclude_modules>], where<model_string>matches an entry in--model_paths/--model_id_with_origin_paths,methodis a registered method (e.g.bitsandbytes_nf4), andexclude_modulesoptionally lists layers kept in full precision.
- Training Basic Configuration
--learning_rate: Learning rate.--num_epochs: Number of epochs.--trainable_models: Trainable models, e.g.,dit,vae,text_encoder.--find_unused_parameters: Whether there are unused parameters in DDP training. Some models contain redundant parameters that do not participate in gradient calculation, and this setting needs to be enabled to avoid errors in multi-GPU training.--weight_decay: Weight decay size, see torch.optim.AdamW.--task: Training task, default issft. Some models support more training modes, please refer to the documentation of each specific model.
- Output Configuration
--output_path: Model saving path.--remove_prefix_in_ckpt: Remove prefix in the state dict of the model file.--save_steps: Interval of training steps to save the model. If this parameter is left blank, the model is saved once per epoch.
- LoRA Configuration
--lora_base_model: Which model to add LoRA to.--lora_target_modules: Which layers to add LoRA to.--lora_rank: Rank of LoRA.--lora_checkpoint: Path of the LoRA checkpoint. If this path is provided, LoRA will be loaded from this checkpoint.--preset_lora_path: Preset LoRA checkpoint path. If this path is provided, this LoRA will be loaded in the form of being merged into the base model. This parameter is used for LoRA differential training.--preset_lora_model: Model that the preset LoRA is merged into, e.g.,dit.
- Gradient Configuration
--use_gradient_checkpointing: Whether to enable gradient checkpointing.--use_gradient_checkpointing_offload: Whether to offload gradient checkpointing to memory.--gradient_accumulation_steps: Number of gradient accumulation steps.
- Video Width/Height Configuration
--height: Height of the video. Leaveheightandwidthblank to enable dynamic resolution.--width: Width of the video. Leaveheightandwidthblank to enable dynamic resolution.--max_pixels: Maximum pixel area of video frames. When dynamic resolution is enabled, video frames with resolution larger than this value will be downscaled, and video frames with resolution smaller than this value will remain unchanged.--num_frames: Number of frames in the video.
- Dataset Basic Configuration
- LTX-2 Series Specific Parameters
--tokenizer_path: Path of the tokenizer, applicable to text-to-video models, leave blank to automatically download from remote.--frame_rate: frame rate of the training videos.
We have built a sample video dataset for your testing. You can download this dataset with the following command:
modelscope download --dataset DiffSynth-Studio/diffsynth_example_dataset --local_dir ./data/diffsynth_example_datasetWe have written recommended training scripts for each model, please refer to the table in the "Model Overview" section above. For how to write model training scripts, please refer to Model Training; for more advanced training algorithms, please refer to Training Framework Detailed Explanation.
En Model Details LingBot Video (`docs/en/Model_Details/LingBot-Video.md`)
LingBot-Video
LingBot-Video is a flow-matching video generation model developed by the LingBot team; a single model handles text-to-video, image-to-video and text-to-image tasks.
Huge thanks to NancyFyong for the outstanding contribution to the integration of this model!
Installation
Before performing model inference and training, please install DiffSynth-Studio first.
git clone https://github.com/modelscope/DiffSynth-Studio.git
cd DiffSynth-Studio
pip install -e .For more information on installation, please refer to Setup Dependencies.
Quick Start
Running the following code will load the Robbyant/lingbot-video-dense-1.3b model for inference. VRAM management is enabled, the framework automatically controls parameter loading based on available VRAM, requiring a minimum of 6GB VRAM.
/* Detailed source-code truncated for AI context efficiency. */Model Overview
| Model ID | Inference | Low VRAM Inference | Full Training | Full Training Validation | LoRA Training | LoRA Training Validation |
|---|---|---|---|---|---|---|
| Robbyant/lingbot-video-dense-1.3b: T2V | code | code | code | code | code | code |
| Robbyant/lingbot-video-dense-1.3b: TI2V | code | code | code | code | code | code |
| Robbyant/lingbot-video-dense-1.3b: T2I | code | code | - | - | - | - |
| Robbyant/lingbot-video-moe-30b-a3b: T2V | code | code | code | code | code | code |
| Robbyant/lingbot-video-moe-30b-a3b: TI2V | code | code | code | code | code | code |
| Robbyant/lingbot-video-moe-30b-a3b: T2I | code | code | - | - | - | - |
| Robbyant/lingbot-video-moe-30b-a3b: T2V + Refinement | code | code | - | - | - | - |
| Robbyant/lingbot-video-moe-30b-a3b: TI2V + Refinement | code | code | - | - | - | - |
Model Inference
The model is loaded via LingBotVideoPipeline.from_pretrained, see Loading Models for details.
The input parameters for LingBotVideoPipeline inference include:
prompt: Structured-JSON caption (dict) or a plain string describing the content. LingBot-Video is trained on structured captions; the pipeline normalises adictautomatically. Released structured captions ship in the example dataset (see Prompt rewriting below).negative_prompt: Negative prompt describing content that should not appear.pipe.default_negative_promptships the official T2V/V2V/TI2V negative prompt;pipe.default_negative_prompt_imageis the T2I variant with temporal terms removed.input_image: First-frame PIL image for image-to-video (TI2V). The frame is VAE-encoded to a clean latent pinned into the first temporal slot after every scheduler step, so the model only generates the frames that follow. LeaveNonefor T2V / V2V / T2I.input_video: Input video (a list of frames or aVideoData) for video-to-video generation, used together withdenoising_strength.denoising_strength: Denoising strength in[0, 1], default1.0. Lower values keep more of the input video structure. Only effective wheninput_videois provided.height: Video / image height, default480. Must be a multiple of 16.width: Video / image width, default480. Must be a multiple of 16.num_frames: Number of frames, default81. Must satisfy4k+1(the VAE compresses time by 4×). Usenum_frames=1for text-to-image.cfg_scale: Classifier-free guidance scale, default3.0.num_inference_steps: Number of inference steps, default40.sigma_shift: Flow-matching timestep shift, default3.0.t_thresh: Refinement start sigma, defaultNone(plain generation). When set, the schedule is truncated so that sampling starts atsigma=t_threshandinput_videois noised to exactly that level. Only meaningful together withinput_video; TI2V additionally re-pins the clean first-frame latent after every step. The official refiner setting is0.85.sigma_tail_steps: Number of extra low-noise steps appended to the tail of the refinement schedule, default2. Only effective whent_threshis set.seed: Random seed. Default isNone, meaning completely random.rand_device: Device for generating the initial noise, default"cpu".progress_bar_cmd: Progress bar, defaulttqdm. Can be disabled by setting tolambda x: x.
If VRAM is insufficient, please enable VRAM Management. We provide recommended low-VRAM configurations for each task in the example code, see the table in the "Model Overview" section above.
Two-stage refinement
The MoE refiner performs a short second pass at a higher resolution: the official setup generates at 480×832 with 40 steps, then refines at 1088×1920 with 8 steps. Load the pipeline with the refiner/ shards instead of transformer/, feed the base clip back in through input_video at the higher resolution, and set t_thresh:
input_video = VideoData("video_base.mp4", height=1088, width=1920)
video = pipe(
prompt=caption,
negative_prompt=pipe.default_negative_prompt,
input_video=input_video,
height=1088, width=1920, num_frames=81,
num_inference_steps=8, cfg_scale=3.0,
t_thresh=0.85, sigma_tail_steps=2,
seed=0,
)The upscaled clip is VAE-encoded and noised back to sigma=t_thresh, so the pass keeps the structure of the base clip and regenerates detail at the target resolution. Pass the same caption as the base pass and keep the same aspect ratio. The refinement resolution dominates the cost — at 1088×1920 the sequence is ~5× longer than at 480×832 — so run this pass with VRAM management enabled.
Prompt rewriting
LingBot-Video is trained on structured-JSON captions, not free-form prose. Feeding a flat sentence is out-of-distribution and visibly degrades quality. The pipeline accepts a caption as a `dic